Miscellany

Jai Ganesh · 2024-09-14 17:01:52

2301) Longitude

Gist

Longitude measures distance east or west of the prime meridian. Lines of longitude, also called meridians, are imaginary lines that divide the Earth. They run north to south from pole to pole, but they measure the distance east or west. Longitude is measured in degrees, minutes, and seconds.

Summary

Longitude is a geographic coordinate that specifies the east–west position of a point on the surface of the Earth, or another celestial body. It is an angular measurement, usually expressed in degrees and denoted by the Greek letter lambda. Meridians are imaginary semicircular lines running from pole to pole that connect points with the same longitude. The prime meridian defines 0° longitude; by convention the International Reference Meridian for the Earth passes near the Royal Observatory in Greenwich, south-east London on the island of Great Britain. Positive longitudes are east of the prime meridian, and negative ones are west.

Because of the Earth's rotation, there is a close connection between longitude and time measurement. Scientifically precise local time varies with longitude: a difference of 15° longitude corresponds to a one-hour difference in local time, due to the differing position in relation to the Sun. Comparing local time to an absolute measure of time allows longitude to be determined. Depending on the era, the absolute time might be obtained from a celestial event visible from both locations, such as a lunar eclipse, or from a time signal transmitted by telegraph or radio. The principle is straightforward, but in practice finding a reliable method of determining longitude took centuries and required the effort of some of the greatest scientific minds.

A location's north–south position along a meridian is given by its latitude, which is approximately the angle between the equatorial plane and the normal from the ground at that location.

Longitude is generally given using the geodetic normal or the gravity direction. The astronomical longitude can differ slightly from the ordinary longitude because of vertical deflection, small variations in Earth's gravitational field.

Details

Lines of longitude, also called meridians, are imaginary lines that divide the Earth. They run north to south from pole to pole, but they measure the distance east or west. Longitude is measured in degrees, minutes, and seconds. Although latitude lines are always equally spaced, longitude lines are furthest from each other at the equator and meet at the poles.

Unlike the equator (which is halfway between the Earth’s north and south poles), the prime meridian is an arbitrary line. In 1884, representatives at the International Meridian Conference in Washington, D.C., met to define the meridian that would represent 0 degrees longitude. For its location, the conference chose a line that ran through the telescope at the Royal Observatory in Greenwich, England. At the time, many nautical charts and time zones already used Greenwich as the starting point, so keeping this location made sense. But, if you go to Greenwich with your GPS receiver, you’ll need to walk 102 meters (334 feet) east of the prime meridian markers before your GPS shows 0 degrees longitude. In the 19th century, scientists did not take into account local variations in gravity or the slightly squished shape of the Earth when they determined the location of the prime meridian. Satellite technology, however, allows scientists to more precisely plot meridians so that they are straight lines running north and south, unaffected by local gravity changes. In the 1980s, the International Reference Meridian (IRM) was established as the precise location of 0 degrees longitude. Unlike the prime meridian, the IRM is not a fixed location, but will continue to move as the Earth’s surface shifts.

Lines of longitude, also called meridians, are imaginary lines that divide the Earth. They run north to south from pole to pole, but they measure the distance east or west.

The prime meridian, which runs through Greenwich, England, has a longitude of 0 degrees. It divides the Earth into the eastern and western hemispheres. The antimeridian is on the opposite side of the Earth, at 180 degrees longitude. Though the antimeridian is the basis for the international date line, actual date and time zone boundaries are dependent on local laws. The international date line zigzags around borders near the antimeridian.

Like latitude, longitude is measured in degrees, minutes, and seconds. Although latitude lines are always equally spaced, longitude lines are furthest from each other at the equator and meet at the poles. At the equator, longitude lines are the same distance apart as latitude lines — one degree covers about 111 kilometers (69 miles). But, by 60 degrees north or south, that distance is down to 56 kilometers (35 miles). By 90 degrees north or south (at the poles), it reaches zero.

Navigators and mariners have been able to measure latitude with basic tools for thousands of years. Longitude, however, required more advanced tools and calculations. Starting in the 16th century, European governments began offering huge rewards if anyone could solve “the longitude problem.” Several methods were tried, but the best and simplest way to measure longitude from a ship was with an accurate clock.

A navigator would compare the time at local noon (when the sun is at its highest point in the sky) to an onboard clock that was set to Greenwich Mean Time (the time at the prime meridian). Each hour of difference between local noon and the time in Greenwich equals 15 degrees of longitude. Why? Because the Earth rotates 360 degrees in 24 hours, or 15 degrees per hour. If the sun’s position tells the navigator it’s local noon, and the clock says back in Greenwich, England, it’s 2 p.m., the two-hour difference means the ship’s longitude is 30 degrees west.

But aboard a swaying ship in varying temperatures and salty air, even the most accurate clocks of the age did a poor job of keeping time. It wasn’t until marine chronometers were invented in the 18th century that longitude could be accurately measured at sea.

Accurate clocks are still critical to determining longitude, but now they’re found in GPS satellites and stations. Each GPS satellite is equipped with one or more atomic clocks that provide incredibly precise time measurements, accurate to within 40 nanoseconds (or 40 billionths of a second). The satellites broadcast radio signals with precise timestamps. The radio signals travel at a constant speed (the speed of light), so we can easily calculate the distance between a satellite and GPS receiver if we know precisely how long it took for the signal to travel between them.

Additional Information

Longitude is the measurement east or west of the prime meridian. Longitude is measured by imaginary lines that run around Earth vertically (up and down) and meet at the North and South Poles. These lines are known as meridians. Each meridian measures one arc degree of longitude. The distance around Earth measures 360 degrees.

The meridian that runs through Greenwich, England, is internationally accepted as the line of 0 degrees longitude, or prime meridian. The antimeridian is halfway around the world, at 180 degrees. It is the basis for the International Date Line.

Half of the world, the Eastern Hemisphere, is measured in degrees east of the prime meridian. The other half, the Western Hemisphere, in degrees west of the prime meridian.

Degrees of longitude are divided into 60 minutes. Each minute of longitude can be further divided into 60 seconds. For example, the longitude of Paris, France, is 2° 29' E (2 degrees, 29 minutes east). The longitude for Brasilia, Brazil, is 47° 55' W (47 degrees, 55 minutes west).

A degree of longitude is about 111 kilometers (69 miles) at its widest. The widest areas of longitude are near the Equator, where Earth bulges out. Because of Earth's curvature, the actual distance of a degrees, minutes, and seconds of longitude depends on its distance from the Equator. The greater the distance, the shorter the length between meridians. All meridians meet at the North and South Poles.

Longitude is related to latitude, the measurement of distance north or south of the Equator. Lines of latitude are called parallels. Maps are often marked with parallels and meridians, creating a grid. The point in the grid where parallels and meridians intersect is called a coordinate. Coordinates can be used to locate any point on Earth.

Knowing the exact coordinates of a site (degrees, minutes, and seconds of longitude and latitude) is valuable for military, engineering, and rescue operations. Coordinates can give military leaders the location of weapons or enemy troops. Coordinates help engineers plan the best spot for a building, bridge, well, or other structure. Coordinates help airplane pilots land planes or drop aid packages in specific locations.

Jai Ganesh · 2024-09-15 00:02:34

2302) Fertilizer

Gist

A fertilizer is a natural or artificial substance containing the chemical elements that improve growth and productiveness of plants. Fertilizers enhance the natural fertility of the soil or replace chemical elements taken from the soil by previous crops.

Summary

A fertilizer (American English) or fertiliser (British English) is any material of natural or synthetic origin that is applied to soil or to plant tissues to supply plant nutrients. Fertilizers may be distinct from liming materials or other non-nutrient soil amendments. Many sources of fertilizer exist, both natural and industrially produced. For most modern agricultural practices, fertilization focuses on three main macro nutrients: nitrogen (N), phosphorus (P), and potassium (K) with occasional addition of supplements like rock flour for micronutrients. Farmers apply these fertilizers in a variety of ways: through dry or pelletized or liquid application processes, using large agricultural equipment, or hand-tool methods.

Historically, fertilization came from natural or organic sources: compost, animal manure, human manure, harvested minerals, crop rotations, and byproducts of human-nature industries (e.g. fish processing waste, or bloodmeal from animal slaughter). However, starting in the 19th century, after innovations in plant nutrition, an agricultural industry developed around synthetically created fertilizers. This transition was important in transforming the global food system, allowing for larger-scale industrial agriculture with large crop yields.

Nitrogen-fixing chemical processes, such as the Haber process invented at the beginning of the 20th century, and amplified by production capacity created during World War II, led to a boom in using nitrogen fertilizers. In the latter half of the 20th century, increased use of nitrogen fertilizers (800% increase between 1961 and 2019) has been a crucial component of the increased productivity of conventional food systems (more than 30% per capita) as part of the so-called "Green Revolution".

The use of artificial and industrially-applied fertilizers has caused environmental consequences such as water pollution and eutrophication due to nutritional runoff; carbon and other emissions from fertilizer production and mining; and contamination and pollution of soil. Various sustainable agriculture practices can be implemented to reduce the adverse environmental effects of fertilizer and pesticide use and environmental damage caused by industrial agriculture.

Details

Soil fertility is the quality of a soil that enables it to provide compounds in adequate amounts and proper balance to promote growth of plants when other factors (such as light, moisture, temperature, and soil structure) are favourable. Where fertility of the soil is not good, natural or manufactured materials may be added to supply the needed plant nutrients. These are called fertilizers, although the term is generally applied to largely inorganic materials other than lime or gypsum.

Essential plant nutrients

In total, plants need at least 16 elements, of which the most important are carbon, hydrogen, oxygen, nitrogen, phosphorus, sulfur, potassium, calcium, and magnesium. Plants obtain carbon from the atmosphere and hydrogen and oxygen from water; other nutrients are taken up from the soil. Although plants contain sodium, iodine, and cobalt, these are apparently not essential. This is also true of silicon and aluminum.

Overall chemical analyses indicate that the total supply of nutrients in soils is usually high in comparison with the requirements of crop plants. Much of this potential supply, however, is bound tightly in forms that are not released to crops fast enough to give satisfactory growth. Because of this, the farmer is interested in measuring the available nutrient supply as contrasted to the total nutrient supply. When the available supply of a given nutrient becomes depleted, its absence becomes a limiting factor in plant growth. Excessive quantities of some nutrients may cause a decrease in yield, however.

Determining nutrient needs

Determination of a crop’s nutrient needs is an essential aspect of fertilizer technology. The appearance of a growing crop may indicate a need for fertilizer, though in some plants the need for more or different nutrients may not be easily observable. If such a problem exists, its nature must be diagnosed, the degree of deficiency must be determined, and the amount and kind of fertilizer needed for a given yield must be found. There is no substitute for detailed examination of plants and soil conditions in the field, followed by simple fertilizer tests, quick tests of plant tissues, and analysis of soils and plants.

Sometimes plants show symptoms of poor nutrition. Chlorosis (general yellow or pale green colour), for example, indicates lack of sulfur and nitrogen. Iron deficiency produces white or pale yellow tissue. Symptoms can be misinterpreted, however. Plant disease can produce appearances resembling mineral deficiency, as can various organisms. Drought or improper cultivation or fertilizer application each may create deficiency symptoms.

After field diagnosis, the conclusions may be confirmed by experiments in a greenhouse or by making strip tests in the field. In strip tests, the fertilizer elements suspected of being deficient are added, singly or in combination, and the resulting plant growth observed. Next it is necessary to determine the extent of the deficiency.

An experiment in the field can be conducted by adding nutrients to the crop at various rates. The resulting response of yield in relation to amounts of nutrients supplied will indicate the supplying power of the unfertilized soil in terms of bushels or tons of produce. If the increase in yield is large, this practice will show that the soil has too little of a given nutrient. Such field experiments may not be practical, because they can cost too much in time and money. Soil-testing laboratories are available in most areas; they conduct chemical soil tests to estimate the availability of nutrients. Commercial soil-testing kits give results that may be very inaccurate, depending on techniques and interpretation. Actually, the most accurate system consists of laboratory analysis of the nutrient content of plant parts, such as the leaf. The results, when correlated with yield response to fertilizer application in field experiments, can give the best estimate of deficiency. Further development of remote sensing techniques, such as infrared photography, are under study and may ultimately become the most valuable technique for such estimates.

The economics of fertilizers

The practical goal is to determine how much nutrient material to add. Since the farmer wants to know how much profit to expect when buying fertilizer, the tests are interpreted as an estimation of increased crop production that will result from nutrient additions. The cost of nutrients must be balanced against the value of the crop or even against alternative procedures, such as investing the money in something else with a greater potential return. The law of diminishing returns is well exemplified in fertilizer technology. Past a certain point, equal inputs of chemicals produce less and less yield increase. The goal of the farmer is to use fertilizer in such a way that the most profitable application rate is employed. Ideal fertilizer application also minimizes excess and ill-timed application, which is not only wasteful for the farmer but also harmful to nearby waterways. Unfortunately, water pollution from fertilizer runoff, which has a sphere of impact that extends far beyond the farmer and the fields, is a negative externality that is not accounted for in the costs and prices of the unregulated market.

Fertilizers can aid in making profitable changes in farming. Operators can reduce costs per unit of production and increase the margin of return over total cost by increasing rates of application of fertilizer on principal cash and feed crops. They are then in a position to invest in soil conservation and other improvements that are needed when shifting acreage from surplus crops to other uses.

Synthetic fertilizers

Modern chemical fertilizers include one or more of the three elements that are most important in plant nutrition: nitrogen, phosphorus, and potassium. Of secondary importance are the elements sulfur, magnesium, and calcium.

Most nitrogen fertilizers are obtained from synthetic ammonia; this chemical compound (NH3) is used either as a gas or in a water solution, or it is converted into salts such as ammonium sulfate, ammonium nitrate, and ammonium phosphate, but packinghouse wastes, treated garbage, sewage, and manure are also common sources of it. Because its nitrogen content is high and is readily converted to ammonia in the soil, urea is one of the most concentrated nitrogenous fertilizers. An inexpensive compound, it is incorporated in mixed fertilizers as well as being applied alone to the soil or sprayed on foliage. With formaldehyde it gives methylene-urea fertilizers, which release nitrogen slowly, continuously, and uniformly, a full year’s supply being applied at one time.

Phosphorus fertilizers include calcium phosphate derived from phosphate rock or bones. The more soluble superphosphate and triple superphosphate preparations are obtained by the treatment of calcium phosphate with sulfuric and phosphoric acid, respectively. Potassium fertilizers, namely potassium chloride and potassium sulfate, are mined from potash deposits. Of commercially produced potassium compounds, almost 95 percent of them are used in agriculture as fertilizer.

Mixed fertilizers contain more than one of the three major nutrients—nitrogen, phosphorus, and potassium. Fertilizer grade is a conventional expression that indicates the percentage of plant nutrients in a fertilizer; thus, a 10–20–10 grade contains 10 percent nitrogen, 20 percent phosphoric oxide, and 10 percent potash. Mixed fertilizers can be formulated in hundreds of ways.

Organic fertilizers and practices

The use of manure and compost as fertilizers is probably almost as old as agriculture. Many traditional farming systems still rely on these sustainable fertilizers, and their use is vital to the productivity of certified organic farms, in which synthetic fertilizers are not permitted.

Farm manure

Among sources of organic matter and plant nutrients, farm manure has been of major importance. Manure is understood to mean the refuse from stables and barnyards, including both excreta and straw or other bedding material. Large amounts of manure are produced by livestock; such manure has value in maintaining and improving soil because of the plant nutrients, humus, and organic substances contained in it.

Due to the potential for harbouring human pathogens, the USDA National Organic Standards mandate that raw manure must be applied no later than 90 or 120 days before harvest, depending on whether the harvested part of the crop is in contact with the ground. Composted manure that has been turned five times in 15 days and reached temperatures between 55 and 77.2 °C (131 and 171 °F) has no restrictions on application times. As manure must be managed carefully in order to derive the most benefit from it, some farmers may be unwilling to expend the necessary time and effort. Manure must be carefully stored to minimize loss of nutrients, particularly nitrogen. It must be applied to the right kind of crop at the proper time. Also, additional fertilizer may be needed, such as phosphoric oxide, in order to gain full value of the nitrogen and potash that are contained in manure.

Manure is fertilizer graded as approximately 0.5–0.25–0.5 (percentages of nitrogen, phosphoric oxide, and potash), with at least two-thirds of the nitrogen in slow-acting forms. Given that these nutrients are mostly in an unmineralized form that cannot be taken up by plants, soil microbes are needed to break down organic matter and transform nutrients into a bioavailable “mineralized” state. In comparison, synthetic fertilizers are already in mineralized form and can be taken up by plants directly. On properly tilled soils, the returns from synthetic fertilizer usually will be greater than from an equivalent amount of manure. However, manure provides many indirect benefits. It supplies humus, which improves the soil’s physical character by increasing its capacity to absorb and store water, by enhancement of aeration, and by favouring the activities of lower organisms. Manure incorporated into the topsoil will help prevent erosion from heavy rain and slow down evaporation of water from the surface. In effect, the value of manure as a mulching material may be greater than is its value as a source of essential plant nutrients.

Green manuring

In reasonably humid areas the practice of green manuring can improve yield and soil qualities. A green manure crop is grown and plowed under for its beneficial effects, although during its growth it may be grazed. These green manure crops are usually annuals, either grasses or legumes, whose roots bear nodule bacteria capable of fixing atmospheric nitrogen. Among the advantages of green manure crops are the addition of nitrogen to the soil, an increase in general fertility, a reduction of erosion, an improvement of physical condition, and a reduction of nutrient loss from leaching. Disadvantages include the chance of not obtaining satisfactory growth; the possibility that the cost of growing the manure crop may exceed the cost of applying commercial nitrogen; possible increases in disease, insect pests, and nematodes (parasitic worms); and possible exhaustion of soil moisture by the crop.

Green manure crops are usually planted in the fall and turned under in the spring before the summer crop is sown. Their value as a source of nitrogen, particularly that of the legumes, is unquestioned for certain crops such as potatoes, cotton, and corn (maize); for other crops, such as peanuts (groundnuts; themselves legumes), the practice is questionable.

Compost

Compost is used in agriculture and gardening primarily as a soil amendment rather than as fertilizer, because it has a low content of plant nutrients. It may be incorporated into the soil or mulched on the surface. Heavy rates of application are common.

Compost is basically a mass of rotted organic matter made from waste plant residues. Addition of nitrogen during decomposition is usually advisable. The result is a crumbly material that when added to soil does not compete with the crop for nitrogen. When properly prepared, it is free of obnoxious odours. Composts commonly contain about 2 percent nitrogen, 0.5 to 1 percent phosphorus, and about 2 percent potassium. The nitrogen of compost becomes available slowly and never approaches that available from inorganic sources. This slow release of nitrogen reduces leaching and extends availability over the whole growing season. Composts are essentially fertilizers with low nutrient content, which explains why large amounts are applied. The maximum benefits of composts on soil structure (better aggregation, pore spacing, and water storage) and on crop yield usually occur after several years of use.

In practical farming, the use of composted plant residues must be compared with the use of fresh residues. More beneficial soil effects usually accrue with less labour by simply turning under fresh residues; also, since one-half the organic matter is lost in composting, fresh residues applied at the same rate will cover twice the area that composted residues would cover. In areas where commercial fertilizers are expensive, labour is cheap, and implements are simple, however, composting meets the need and is a logical practice.

Sewage sludge, the solid material remaining from the treatment of sewage, is not permitted in certified organic farming, though it is used in other, nonorganic settings. After suitable processing, it is sold as fertilizer and as a soil amendment for use on lawns, in parks, and on golf courses. Use of human biosolids in agriculture is controversial, as there are concerns that even treated sewage may harbour harmful bacteria, viruses, pharmaceutical residues, and heavy metals.

Liming

Liming to reduce soil acidity is practiced extensively in humid areas where rainfall leaches calcium and magnesium from the soil, thus creating an acid condition. Calcium and magnesium are major plant nutrients supplied by liming materials. Ground limestone is widely used for this purpose; its active agent, calcium carbonate, reacts with the soil to reduce its acidity. The calcium is then available for plant use. The typical limestones, especially dolomitic, contain magnesium carbonate as well, thus also supplying magnesium to the plant.

Marl and chalk are soft impure forms of limestone and are sometimes used as liming materials, as are oyster shells. Calcium sulfate (gypsum) and calcium chloride, however, are unsuitable for liming, for, although their calcium is readily soluble, they leave behind a residue that is harmful. Organic standards by the European Union and the U.S. Food and Drug Administration restrict certain liming agents; burnt lime and hydrated lime are not permitted for certified organic farms in the U.S., for example.

Lime is applied by mixing it uniformly with the surface layer of the soil. It may be applied at any time of the year on land plowed for spring crops or winter grain or on permanent pasture. After application, plowing, disking, or harrowing will mix it with the soil. Such tillage is usually necessary, because calcium migrates slowly downward in most soils. Lime is usually applied by trucks specially equipped and owned by custom operators.

Methods of application

Fertilizers may be added to soil in solid, liquid, or gaseous forms, the choice depending on many factors. Generally, the farmer tries to obtain satisfactory yield at minimum cost in money and labour.

Manure can be applied as a liquid or a solid. When accumulated as a liquid from livestock areas, it may be stored in tanks until needed and then pumped into a distributing machine or into a sprinkler irrigation system. The method reduces labour, but the noxious odours are objectionable. The solid-manure spreader, which can also be used for compost, conveys the material to the field, shreds it, and spreads it uniformly over the land. The process can be carried out during convenient times, including winter, but rarely when the crop is growing.

Application of granulated or pelleted solid fertilizer has been aided by improved equipment design. Such devices, depending on design, can deposit fertilizer at the time of planting, side-dress a growing crop, or broadcast the material. Solid-fertilizer distributors have a wide hopper with holes in the bottom; distribution is effected by various means, such as rollers, agitators, or endless chains traversing the hopper bottom. Broadcast distributors have a tub-shaped hopper from which the material falls onto revolving disks that distribute it in a broad swath. Fertilizer attachments are available for most tractor-mounted planters and cultivators and for grain drills and some types of plows. They deposit fertilizer with the seed when planted, without damage to the seed, yet the nutrient is readily available during early growth. Placement of the fertilizer varies according to the types of crops; some crops require banding above the seed, while others are more successful when the fertilizer band is below the seed.

The use of liquid and ammonia fertilizers is growing, particularly of anhydrous ammonia, which is handled as a liquid under pressure but changes to gas when released to atmospheric pressure. Anhydrous ammonia, however, is highly corrosive, inflammable, and rather dangerous if not handled properly; thus, application equipment is specialized. Typically, the applicator is a chisel-shaped blade with a pipe mounted on its rear side to conduct the ammonia 13 to 15 cm (5 to 6 inches) below the soil surface. Pipes are fed from a pressure tank mounted above. Mixed liquid fertilizers containing nitrogen, phosphorus, and potassium may be applied directly to the soil surface or as a foliar spray by field sprayers where close-growing crops are raised. Large areas can be covered rapidly by use of aircraft, which can distribute both liquid and dry fertilizer.

Additional Information

Fertilizers have played an essential role in feeding a growing global population. It's estimated that just under half of the people alive today are dependent on synthetic fertilizers.

They can bring environmental benefits too: fertilizers can increase crop yields. By increasing crop yields we can reduce the amount of land we use for agriculture.

But they also create environmental pollution. Many countries overapply fertilizers, leading to the runoff of nutrients into water systems and ecosystems.

A problem we need to tackle is using fertilizers efficiently: yielding its benefits to feed a growing population while reducing the environmental damage that they cause.

Organic and Mineral Fertilizer: Differences and Similarities.

Fertilizers are materials that are applied to soils, or directly to plants, for their ability to supply the essential nutrients needed by crops to grow and improve soil fertility. They are used to increase crop yield and/or quality, as well as to sustain soils’ ability to support future crop production.

Mineral fertilizers are produced from materials mined from naturally occurring nutrient deposits, or from the fixation of nitrogen from the atmosphere into plant-available forms. Mineral fertilizers generally contain high concentrations of a single, or two or three, plant nutrients.

Organic fertilizers are derived from plant matter, animal excreta, sewage and food waste, generally in the form of animal manure, green manure and biosolids. Organic fertilizers provide essential nutrients needed by crops, generally containing a wide variety in low concentrations. They also play an important role in improving soil health.

Organo-mineral fertilizers combine dried organic and mineral fertilizers to provide balanced nutrients along with soil health improvements in a long-lasting, easy-to transport and store form.

They play a key role in reducing micronutrient deficiencies in people:

The fertilizer fortification of staple food crops with micronutrients (also known as agronomic biofortification) has alleviated deficiencies in zinc, selenium and iodine in communities around the world.

The fertilizer industry supports policies that link agriculture, nutrition and health, and the use of micronutrients where they are needed most.

Fertilizing-plants-0440837aeee64749832645ba62572f95.jpg

Jai Ganesh · 2024-09-16 00:04:36

2303) Siren

Gist

A siren is a warning device which makes a long, loud noise. Most fire engines, ambulances, and police cars have sirens.

It is a device for making a loud warning noise.

Summary

A siren is noisemaking device producing a piercing sound of definite pitch. Used as a warning signal, it was invented in the late 18th century by the Scottish natural philosopher John Robison. The name was given it by the French engineer Charles Cagniard de La Tour, who devised an acoustical instrument of the type in 1819. A disk with evenly spaced holes around its edge is rotated at high speed, interrupting at regular intervals a jet of air directed at the holes. The resulting regular pulsations cause a sound wave in the surrounding air. The siren is thus classified as a free aerophone. The sound-wave frequency of its pitch equals the number of air puffs (or holes times number of revolutions) per second. The strident sound results from the high number of overtones (harmonics) present.

Details

A siren is a loud noise-making device. Civil defense sirens are mounted in fixed locations and used to warn of natural disasters or attacks. Sirens are used on emergency service vehicles such as ambulances, police cars, and fire engines. There are two general types: mechanical and electronic.

Many fire sirens (used for summoning volunteer firefighters) serve double duty as tornado or civil defense sirens, alerting an entire community of impending danger. Most fire sirens are either mounted on the roof of a fire station or on a pole next to the fire station. Fire sirens can also be mounted on or near government buildings, on tall structures such as water towers, as well as in systems where several sirens are distributed around a town for better sound coverage. Most fire sirens are single tone and mechanically driven by electric motors with a rotor attached to the shaft. Some newer sirens are electronically driven speakers.

Fire sirens are often called fire whistles, fire alarms, or fire horns. Although there is no standard signaling of fire sirens, some utilize codes to inform firefighters of the location of the fire. Civil defense sirens also used as fire sirens often can produce an alternating "hi-lo" signal (similar to emergency vehicles in many European countries) as the fire signal, or attack (slow wail), typically 3x, as to not confuse the public with the standard civil defense signals of alert (steady tone) and fast wail (fast wavering tone). Fire sirens are often tested once a day at noon and are also called noon sirens or noon whistles.

The first emergency vehicles relied on a bell. In the 1970s, they switched to a duotone airhorn, which was itself overtaken in the 1980s by an electronic wail.

History

Some time before 1799, the siren was invented by the Scottish natural philosopher John Robison. Robison's sirens were used as musical instruments; specifically, they powered some of the pipes in an organ. Robison's siren consisted of a stopcock that opened and closed a pneumatic tube. The stopcock was apparently driven by the rotation of a wheel.

In 1819, an improved siren was developed and named by Baron Charles Cagniard de la Tour. De la Tour's siren consisted of two perforated disks that were mounted coaxially at the outlet of a pneumatic tube. One disk was stationary, while the other disk rotated. The rotating disk periodically interrupted the flow of air through the fixed disk, producing a tone. De la Tour's siren could produce sound under water, suggesting a link with the sirens of Greek mythology; hence the name he gave to the instrument.

Instead of disks, most modern mechanical sirens use two concentric cylinders, which have slots parallel to their length. The inner cylinder rotates while the outer one remains stationary. As air under pressure flows out of the slots of the inner cylinder and then escapes through the slots of the outer cylinder, the flow is periodically interrupted, creating a tone. The earliest such sirens were developed during 1877–1880 by James Douglass and George Slight (1859–1934) of Trinity House; the final version was first installed in 1887 at the Ailsa Craig lighthouse in Scotland's Firth of Clyde. When commercial electric power became available, sirens were no longer driven by external sources of compressed air, but by electric motors, which generated the necessary flow of air via a simple centrifugal fan, which was incorporated into the siren's inner cylinder.

To direct a siren's sound and to maximize its power output, a siren is often fitted with a horn, which transforms the high-pressure sound waves in the siren to lower-pressure sound waves in the open air.

The earliest way of summoning volunteer firemen to a fire was by ringing of a bell, either mounted atop the fire station, or in the belfry of a local church. As electricity became available, the first fire sirens were manufactured. In 1886 French electrical engineer Gustave Trouvé developed a siren to announce the silent arrival of his electric boats. Two early fire siren manufacturers were William A. Box Iron Works, who made the "Denver" sirens as early as 1905, and the Inter-State Machine Company (later the Sterling Siren Fire Alarm Company) who made the ubiquitous Model "M" electric siren, which was the first dual tone siren. The popularity of fire sirens took off by the 1920s, with many manufacturers including the Federal Electric Company and Decot Machine Works creating their own sirens. Since the 1970s, many communities have since deactivated their fire sirens as pagers became available for fire department use. Some sirens still remain as a backup to pager systems.

During the Second World War, the British civil defence used a network of sirens to alert the general population to the imminence of an air raid. A single tone denoted an "all clear". A series of tones denoted an air raid.

Types:

Pneumatic

The pneumatic siren, which is a free aerophone, consists of a rotating disk with holes in it (called a chopper, siren disk or rotor), such that the material between the holes interrupts a flow of air from fixed holes on the outside of the unit (called a stator). As the holes in the rotating disk alternately prevent and allow air to flow it results in alternating compressed and rarefied air pressure, i.e. sound. Such sirens can consume large amounts of energy. To reduce the energy consumption without losing sound volume, some designs of pneumatic sirens are boosted by forcing compressed air from a tank that can be refilled by a low powered compressor through the siren disk.

In United States English language usage, vehicular pneumatic sirens are sometimes referred to as mechanical or coaster sirens, to differentiate them from electronic devices. Mechanical sirens driven by an electric motor are often called "electromechanical". One example is the Q2B siren sold by Federal Signal Corporation. Because of its high current draw (100 amps when power is applied) its application is normally limited to fire apparatus, though it has seen increasing use on type IV ambulances and rescue-squad vehicles. Its distinct tone of urgency, high sound pressure level (123 dB at 10 feet) and square sound waves account for its effectiveness.

In Germany and some other European countries, the pneumatic two-tone (hi-lo) siren consists of two sets of air horns, one high pitched and the other low pitched. An air compressor blows the air into one set of horns, and then it automatically switches to the other set. As this back and forth switching occurs, the sound changes tones. Its sound power varies, but could get as high as approximately 125 dB, depending on the compressor and the horns. Comparing with the mechanical sirens, it uses much less electricity but needs more maintenance.

In a pneumatic siren, the stator is the part which cuts off and reopens air as rotating blades of a chopper move past the port holes of the stator, generating sound. The pitch of the siren's sound is a function of the speed of the rotor and the number of holes in the stator. A siren with only one row of ports is called a single tone siren. A siren with two rows of ports is known as a dual tone siren. By placing a second stator over the main stator and attaching a solenoid to it, one can repeatedly close and open all of the stator ports thus creating a tone called a pulse. If this is done while the siren is wailing (rather than sounding a steady tone) then it is called a pulse wail. By doing this separately over each row of ports on a dual tone siren, one can alternately sound each of the two tones back and forth, creating a tone known as Hi/Lo. If this is done while the siren is wailing, it is called a Hi/Lo wail. This equipment can also do pulse or pulse wail. The ports can be opened and closed to send Morse code. A siren which can do both pulse and Morse code is known as a code siren.

Electronic

Electronic sirens incorporate circuits such as oscillators, modulators, and amplifiers to synthesize a selected siren tone (wail, yelp, pierce/priority/phaser, hi-lo, scan, airhorn, manual, and a few more) which is played through external speakers. It is not unusual, especially in the case of modern fire engines, to see an emergency vehicle equipped with both types of sirens. Often, police sirens also use the interval of a tritone to help draw attention. The first electronic siren that mimicked the sound of a mechanical siren was invented in 1965 by Motorola employees Ronald H. Chapman and Charles W. Stephens.

Other types

Steam whistles were also used as a warning device if a supply of steam was present, such as a sawmill or factory. These were common before fire sirens became widely available, particularly in the former Soviet Union. Fire horns, large compressed air horns, also were and still are used as an alternative to a fire siren. Many fire horn systems were wired to fire pull boxes that were located around a town, and this would "blast out" a code in respect to that box's location. For example, pull box number 233, when pulled, would trigger the fire horn to sound two blasts, followed by a pause, followed by three blasts, followed by a pause, followed by three more blasts. In the days before telephones, this was the only way firefighters would know the location of a fire. The coded blasts were usually repeated several times. This technology was also applied to many steam whistles as well. Some fire sirens are fitted with brakes and dampers, enabling them to sound out codes as well. These units tended to be unreliable, and are now uncommon.

Physics of the sound

Mechanical sirens blow air through a slotted disk or rotor. The cyclic waves of air pressure are the physical form of sound. In many sirens, a centrifugal blower and rotor are integrated into a single piece of material, spun by an electric motor.

Electronic sirens are high efficiency loudspeakers, with specialized amplifiers and tone generation. They usually imitate the sounds of mechanical sirens in order to be recognizable as sirens.

To improve the efficiency of the siren, it uses a relatively low frequency, usually several hundred hertz. Lower frequency sound waves go around corners and through holes better.

Sirens often use horns to aim the pressure waves. This uses the siren's energy more efficiently by aiming it. Exponential horns achieve similar efficiencies with less material.

The frequency, i.e. the cycles per second of the sound of a mechanical siren is controlled by the speed of its rotor, and the number of openings. The wailing of a mechanical siren occurs as the rotor speeds and slows. Wailing usually identifies an attack or urgent emergency.

The characteristic timbre or musical quality of a mechanical siren is caused because it is a triangle wave, when graphed as pressure over time. As the openings widen, the emitted pressure increases. As they close, it decreases. So, the characteristic frequency distribution of the sound has harmonics at odd (1, 3, 5...) multiples of the fundamental. The power of the harmonics roll off in an inverse square to their frequency. Distant sirens sound more "mellow" or "warmer" because their harsh high frequencies are absorbed by nearby objects.

Two tone sirens are often designed to emit a minor third, musically considered a "sad" sound. To do this, they have two rotors with different numbers of openings. The upper tone is produced by a rotor with a count of openings divisible by six. The lower tone's rotor has a count of openings divisible by five. Unlike an organ, a mechanical siren's minor third is almost always physical, not tempered. To achieve tempered ratios in a mechanical siren, the rotors must either be geared, run by different motors, or have very large numbers of openings. Electronic sirens can easily produce a tempered minor third.

A mechanical siren that can alternate between its tones uses solenoids to move rotary shutters that cut off the air supply to one rotor, then the other. This is often used to identify a fire warning.

When testing, a frightening sound is not desirable. So, electronic sirens then usually emit musical tones:
Westminster chimes is common. Mechanical sirens sometimes self-test by "growling", i.e. operating at low speeds.

In music
Sirens are also used as musical instruments. They have been prominently featured in works by avant-garde and contemporary classical composers. Examples include Edgard Varèse's compositions Amériques (1918–21, rev. 1927), Hyperprism (1924), and Ionisation (1931); math Avraamov's Symphony of Factory Sirens (1922); George Antheil's Ballet Mécanique (1926); Dimitri Shostakovich's Symphony No. 2 (1927), and Henry Fillmore's "The Klaxon: March of the Automobiles" (1929), which features a klaxophone.

In popular music, sirens have been used in The Chemical Brothers' "Song to the Siren" (1992) and in a CBS News 60 Minutes segment played by percussionist Evelyn Glennie. A variation of a siren, played on a keyboard, are the opening notes of the REO Speedwagon song "Ridin' the Storm Out". Some heavy metal bands also use air raid type siren intros at the beginning of their shows. The opening measure of Money City Maniacs 1998 by Canadian band Sloan uses multiple sirens overlapped.

Jai Ganesh · 2024-09-16 20:11:30

2304) Stainless Steel

Gist

Stainless steel is a corrosion-resistant alloy of iron, chromium and, in some cases, nickel and other metals. Completely and infinitely recyclable, stainless steel is the “green material” par excellence.

Stainless steel is the name of a family of iron-based alloys known for their corrosion and heat resistance. One of the main characteristics of stainless steel is its minimum chromium content of 10.5%, which gives it its superior resistance to corrosion in comparison to other types of steels.

Summary

Stainless steels are a family of ferrous alloys containing less than 1.2% carbon and over 10.5% chromium and are protected by a passive surface layer of chromium and iron oxides and hydroxides that protects them efficiently from corrosion.

Stainless steel is a family of alloy steels containing low carbon steel with a minimum chromium content of 10% or more by weight. The name originates from the fact that stainless steel does not stain, corrode or rust as easily as ordinary steel; however, ‘stain-less’ is not ‘stain-proof’ in all conditions. It is important to select the correct type and grade of stainless steel for a particular application. In many cases, manufacturing rooms, processing lines, equipment and machines will be subject to requirements from authorities, manufacturers or customers.

The addition of chromium gives the steel its unique stainless, corrosion-resistant properties. The chromium, when in contact with oxygen, forms a natural barrier of adherent chromium(III) oxide (Cr2O3), commonly called ‘ceramic,’ which is a ‘passive film’ resistant to further oxidation or rusting. This event is called passivation and is seen in other metals, such as aluminum and silver, but unlike in these metals this passive film is transparent on stainless steel. This invisible, self repairing and relatively inert film is only a few microns thick so the metal stays shiny. If damaged mechanically or chemically, the film is self-healing, meaning the layer quickly reforms, providing that oxygen is present, even if in very small amounts. This protective oxide or ceramic coating is common to most corrosion resistant materials. Similarly, anodizing is an electrolytic passivation process used to increase the thickness of the natural oxide layer on the surface of metals such as aluminum, titanium, and zinc among others. Passivation is not a useful treatment for iron or carbon steel because these metals exfoliate when oxidized, i.e. the iron(III) oxide (rust) flakes off, constantly exposing the underlying metal to corrosion.

The corrosion resistance and other useful properties of stainless steel can be enhanced by increasing the chromium content and the addition of other alloying elements such as molybdenum, nickel and nitrogen. There are more than 60 grades of stainless steel, however, the entire group can be divided into five classes (cast stainless steels, in general, are similar to the equivalent wrought alloys). Each is identified by the alloying elements which affect their microstructure and for which each is named.

Details

Stainless steel, also known as inox, corrosion-resistant steel (CRES), and rustless steel, is an alloy of iron that is resistant to rusting and corrosion. It contains iron with chromium and other elements such as molybdenum, carbon, nickel and nitrogen depending on its specific use and cost. Stainless steel's resistance to corrosion results from the 10.5%, or more, chromium content which forms a passive film that can protect the material and self-heal in the presence of oxygen.

The alloy's properties, such as luster and resistance to corrosion, are useful in many applications. Stainless steel can be rolled into sheets, plates, bars, wire, and tubing. These can be used in cookware, cutlery, surgical instruments, major appliances, vehicles, construction material in large buildings, industrial equipment (e.g., in paper mills, chemical plants, water treatment), and storage tanks and tankers for chemicals and food products. Some grades are also suitable for forging and casting.

The biological cleanability of stainless steel is superior to both aluminium and copper, and comparable to glass. Its cleanability, strength, and corrosion resistance have prompted the use of stainless steel in pharmaceutical and food processing plants.

Different types of stainless steel are labeled with an AISI three-digit number. The ISO 15510 standard lists the chemical compositions of stainless steels of the specifications in existing ISO, ASTM, EN, JIS, and GB standards in a useful interchange table.

Properties:

Corrosion resistance

Although stainless steel does rust, this only affects the outer few layers of atoms, its chromium content shielding deeper layers from oxidation.

The addition of nitrogen also improves resistance to pitting corrosion and increases mechanical strength. Thus, there are numerous grades of stainless steel with varying chromium and molybdenum contents to suit the environment the alloy must endure. Corrosion resistance can be increased further by the following means:

* increasing chromium content to more than 11%
* adding nickel to at least 8%
* adding molybdenum (which also improves resistance to pitting corrosion)

Strength

The most common type of stainless steel, 304, has a tensile yield strength around 210 MPa (30,000 psi) in the annealed condition. It can be strengthened by cold working to a strength of 1,050 MPa (153,000 psi) in the full-hard condition.

The strongest commonly available stainless steels are precipitation hardening alloys such as 17-4 PH and Custom 465. These can be heat treated to have tensile yield strengths up to 1,730 MPa (251,000 psi).

Melting point

Stainless steel is a steel, and as such its melting point is near that of ordinary steel, and much higher than the melting points of aluminium or copper. As with most alloys, the melting point of stainless steel is expressed in the form of a range of temperatures, and not a single temperature. This temperature range goes from 1,400 to 1,530 °C (2,550 to 2,790 °F; 1,670 to 1,800 K; 3,010 to 3,250 °R) depending on the specific consistency of the alloy in question.

Conductivity

Like steel, stainless steels are relatively poor conductors of electricity, with significantly lower electrical conductivities than copper. In particular, the non-electrical contact resistance (ECR) of stainless steel arises as a result of the dense protective oxide layer and limits its functionality in applications as electrical connectors. Copper alloys and nickel-coated connectors tend to exhibit lower ECR values and are preferred materials for such applications. Nevertheless, stainless steel connectors are employed in situations where ECR poses a lower design criteria and corrosion resistance is required, for example in high temperatures and oxidizing environments.

Magnetism

Martensitic, duplex and ferritic stainless steels are magnetic, while austenitic stainless steel is usually non-magnetic. Ferritic steel owes its magnetism to its body-centered cubic crystal structure, in which iron atoms are arranged in cubes (with one iron atom at each corner) and an additional iron atom in the center. This central iron atom is responsible for ferritic steel's magnetic properties. This arrangement also limits the amount of carbon the steel can absorb to around 0.025%. Grades with low coercive field have been developed for electro-valves used in household appliances and for injection systems in internal combustion engines. Some applications require non-magnetic materials, such as magnetic resonance imaging. Austenitic stainless steels, which are usually non-magnetic, can be made slightly magnetic through work hardening. Sometimes, if austenitic steel is bent or cut, magnetism occurs along the edge of the stainless steel because the crystal structure rearranges itself.

Wear

Galling, sometimes called cold welding, is a form of severe adhesive wear, which can occur when two metal surfaces are in relative motion to each other and under heavy pressure. Austenitic stainless steel fasteners are particularly susceptible to thread galling, though other alloys that self-generate a protective oxide surface film, such as aluminum and titanium, are also susceptible. Under high contact-force sliding, this oxide can be deformed, broken, and removed from parts of the component, exposing the bare reactive metal. When the two surfaces are of the same material, these exposed surfaces can easily fuse. Separation of the two surfaces can result in surface tearing and even complete seizure of metal components or fasteners. Galling can be mitigated by the use of dissimilar materials (bronze against stainless steel) or using different stainless steels (martensitic against austenitic). Additionally, threaded joints may be lubricated to provide a film between the two parts and prevent galling. Nitronic 60, made by selective alloying with manganese, silicon, and nitrogen, has demonstrated a reduced tendency to gall.

Density

The density of stainless steel ranges from 7.5 to 8.0 g/{cm}^3 (0.27 to 0.29 lb/cu in) depending on the alloy.

Additional Information

Stainless steel is any one of a family of alloy steels usually containing 10 to 30 percent chromium. In conjunction with low carbon content, chromium imparts remarkable resistance to corrosion and heat. Other elements, such as nickel, molybdenum, titanium, aluminum, niobium, copper, nitrogen, sulfur, phosphorus, or selenium, may be added to increase corrosion resistance to specific environments, enhance oxidation resistance, and impart special characteristics.

Most stainless steels are first melted in electric-arc or basic oxygen furnaces and subsequently refined in another steelmaking vessel, mainly to lower the carbon content. In the argon-oxygen decarburization process, a mixture of oxygen and argon gas is injected into the liquid steel. By varying the ratio of oxygen and argon, it is possible to remove carbon to controlled levels by oxidizing it to carbon monoxide without also oxidizing and losing expensive chromium. Thus, cheaper raw materials, such as high-carbon ferrochromium, may be used in the initial melting operation.

There are more than 100 grades of stainless steel. The majority are classified into five major groups in the family of stainless steels: austenitic, ferritic, martensitic, duplex, and precipitation-hardening. Austenitic steels, which contain 16 to 26 percent chromium and up to 35 percent nickel, usually have the highest corrosion resistance. They are not hardenable by heat treatment and are nonmagnetic. The most common type is the 18/8, or 304, grade, which contains 18 percent chromium and 8 percent nickel. Typical applications include aircraft and the dairy and food-processing industries. Standard ferritic steels contain 10.5 to 27 percent chromium and are nickel-free; because of their low carbon content (less than 0.2 percent), they are not hardenable by heat treatment and have less critical anticorrosion applications, such as architectural and auto trim. Martensitic steels typically contain 11.5 to 18 percent chromium and up to 1.2 percent carbon with nickel sometimes added. They are hardenable by heat treatment, have modest corrosion resistance, and are employed in cutlery, surgical instruments, wrenches, and turbines. Duplex stainless steels are a combination of austenitic and ferritic stainless steels in equal amounts; they contain 21 to 27 percent chromium, 1.35 to 8 percent nickel, 0.05 to 3 percent copper, and 0.05 to 5 percent molybdenum. Duplex stainless steels are stronger and more resistant to corrosion than austenitic and ferritic stainless steels, which makes them useful in storage-tank construction, chemical processing, and containers for transporting chemicals. Precipitation-hardening stainless steel is characterized by its strength, which stems from the addition of aluminum, copper, and niobium to the alloy in amounts less than 0.5 percent of the alloy’s total mass. It is comparable to austenitic stainless steel with respect to its corrosion resistance, and it contains 15 to 17.5 percent chromium, 3 to 5 percent nickel, and 3 to 5 percent copper. Precipitation-hardening stainless steel is used in the construction of long shafts.

Jai Ganesh · 2024-09-17 00:03:00

2305) Earphones / Headphones

Gist

Headphones are also known as earphones or, colloquially, cans. Circumaural (around the ear) and supra-aural (over the ear) headphones use a band over the top of the head to hold the drivers in place. Another type, known as earbuds or earpieces, consists of individual units that plug into the user's ear canal.

They are electroacoustic transducers, which convert an electrical signal to a corresponding sound. Headphones let a single user listen to an audio source privately, in contrast to a loudspeaker, which emits sound into the open air for anyone nearby to hear. Headphones are also known as earphones or, colloquially, cans.

Headphones are a type of hardware output device that can be connected to a computer's line-out or speakers port, as well as wirelessly using Bluetooth. They are also referred to as earbuds. You can watch a movie or listen to audio without bothering anyone nearby by using headphones.

Summary

A headphone is a small loudspeaker (earphone) held over the ear by a band or wire worn on the head. Headphones are commonly employed in situations in which levels of surrounding noise are high, as in an airplane math, or where a user such as a switchboard operator needs to keep the hands free, or where the listener is moving about or wants to listen without disturbing other people. A headphone may be equipped with one earphone or two and may include a miniature microphone, in which case it is called a headset. For listening to stereophonically reproduced sound, stereo headphones may be used, with separate channels of sound being fed to the two earphones.

An earphone is a small loudspeaker held or worn close to the listener’s ear or within the outer ear. Common forms include the hand-held telephone receiver; the headphone, in which one or two earphones are held in place by a band worn over the head; and the plug earphone, which is inserted in the outer opening of the ear. The conversion of electrical to acoustical signals is effected by any of the devices used in larger loudspeakers; the highest fidelity is provided by the so-called dynamic earphone, which ordinarily is made part of a headphone and equipped with a cushion to isolate the ears from other sound sources.

Details

Headphones are a pair of small loudspeaker drivers worn on or around the head over a user's ears. They are electroacoustic transducers, which convert an electrical signal to a corresponding sound. Headphones let a single user listen to an audio source privately, in contrast to a loudspeaker, which emits sound into the open air for anyone nearby to hear. Headphones are also known as earphones or, colloquially, cans. Circumaural (around the ear) and supra-aural (over the ear) headphones use a band over the top of the head to hold the drivers in place. Another type, known as earbuds or earpieces, consists of individual units that plug into the user's ear canal. A third type are bone conduction headphones, which typically wrap around the back of the head and rest in front of the ear canal, leaving the ear canal open. In the context of telecommunication, a headset is a combination of a headphone and microphone.

Headphones connect to a signal source such as an audio amplifier, radio, CD player, portable media player, mobile phone, video game console, or electronic musical instrument, either directly using a cord, or using wireless technology such as Bluetooth, DECT or FM radio. The first headphones were developed in the late 19th century for use by switchboard operators, to keep their hands free. Initially, the audio quality was mediocre and a step forward was the invention of high fidelity headphones.

Headphones exhibit a range of different audio reproduction quality capabilities. Headsets designed for telephone use typically cannot reproduce sound with the high fidelity of expensive units designed for music listening by audiophiles. Headphones that use cables typically have either a 1/4 inch (6.4 mm) or 1/8 inch (3.2 mm) phone jack for plugging the headphones into the audio source. Some headphones are wireless, using Bluetooth connectivity to receive the audio signal by radio waves from source devices like cellphones and digital players. As a result of the Walkman effect, beginning in the 1980s, headphones started to be used in public places such as sidewalks, grocery stores, and public transit. Headphones are also used by people in various professional contexts, such as audio engineers mixing sound for live concerts or sound recordings and DJs, who use headphones to cue up the next song without the audience hearing, aircraft pilots and call center employees. The latter two types of employees use headphones with an integrated microphone.

History

Headphones grew out of the need to free up a person's hands when operating a telephone. By the 1880s, soon after the invention of the telephone, telephone switchboard operators began to use head apparatuses to mount the telephone receiver. The receiver was mounted on the head by a clamp which held it next to the ear. The head mount freed the switchboard operator's hands, so that they could easily connect the wires of the telephone callers and receivers. The head-mounted telephone receiver in the singular form was called a headphone. These head-mounted phone receivers, unlike modern headphones, only had one earpiece.

By the 1890s a listening device with two earpieces was developed by the British company Electrophone. The device created a listening system through the phone lines that allowed the customer to connect into live feeds of performances at theaters and opera houses across London. Subscribers to the service could listen to the performance through a pair of massive earphones that connected below the chin and were held by a long rod.

French engineer Ernest Mercadier in 1891 patented a set of in-ear headphones. He was awarded U.S. Patent No. 454,138 for "improvements in telephone-receivers...which shall be light enough to be carried while in use on the head of the operator." The German company Siemens Brothers at this time was also selling headpieces for telephone operators which had two earpieces, although placed outside the ear. These headpieces by Siemens Brothers looked fairly similar to modern headphones. The majority of headgear used by telephone operators continued to have only one earpiece.

Modern headphones subsequently evolved out of the emerging field of wireless telegraphy, which was the beginning stage of radio broadcasting. Some early wireless telegraph developers chose to use the telephone receiver's speaker as the detector for the electrical signal of the wireless receiving circuit. By 1902 wireless telegraph innovators, such as Lee de Forest, were using two jointly head-mounted telephone receivers to hear the signal of the receiving circuit. The two head-mounted telephone receivers were called in the singular form "head telephones". By 1908 the headpiece began to be written simply as "head phones", and a year later the compound word "headphones" began to be used.

One of the earliest companies to make headphones for wireless operators was the Holtzer-Cabot Company in 1909. They were also makers of head receivers for telephone operators and normal telephone receivers for the home. Another early manufacturer of headphones was Nathaniel Baldwin. He was the first major supplier of headsets to the U.S. Navy. In 1910 he invented a prototype telephone headset due to his inability to hear sermons during Sunday service. He offered it for testing to the navy, which promptly ordered 100 of them because of their good quality. Wireless Specialty Apparatus Co., in partnership with Baldwin Radio Company, set up a manufacturing facility in Utah to fulfill orders.

These early headphones used moving iron drivers, with either single-ended or balanced armatures. The common single-ended type used voice coils wound around the poles of a permanent magnet, which were positioned close to a flexible steel diaphragm. The audio current through the coils varied the magnetic field of the magnet, exerting a varying force on the diaphragm, causing it to vibrate, creating sound waves. The requirement for high sensitivity meant that no damping was used, so the frequency response of the diaphragm had large peaks due to resonance, resulting in poor sound quality. These early models lacked padding, and were often uncomfortable to wear for long periods. Their impedance varied; headphones used in telegraph and telephone work had an impedance of 75 ohms. Those used with early wireless radio had more turns of finer wire to increase sensitivity. Impedance of 1,000 to 2,000 ohms was common, which suited both crystal sets and triode receivers. Some very sensitive headphones, such as those manufactured by Brandes around 1919, were commonly used for early radio work.

In early powered radios, the headphone was part of the vacuum tube's plate circuit and carried dangerous voltages. It was normally connected directly to the positive high voltage battery terminal, and the other battery terminal was securely grounded. The use of bare electrical connections meant that users could be shocked if they touched the bare headphone connections while adjusting an uncomfortable headset.

In 1958, John C. Koss, an audiophile and jazz musician from Milwaukee, produced the first stereo headphones.

Smaller earbud type earpieces, which plugged into the user's ear canal, were first developed for hearing aids. They became widely used with transistor radios, which commercially appeared in 1954 with the introduction of the Regency TR-1. The most popular audio device in history, the transistor radio changed listening habits, allowing people to listen to radio anywhere. The earbud uses either a moving iron driver or a piezoelectric crystal to produce sound. The 3.5 mm radio and phone connector, which is the most commonly used in portable application today, has been used at least since the Sony EFM-117J transistor radio, which was released in 1964. Its popularity was reinforced with its use on the Walkman portable tape player in 1979.

Applications

Headphones may be used with stationary CD and DVD players, home theater, personal computers, or portable devices (e.g., digital audio player/MP3 player, mobile phone), as long as these devices are equipped with a headphone jack. Cordless headphones are not connected to their source by a cable. Instead, they receive a radio or infrared signal encoded using a radio or infrared transmission link, such as FM, Bluetooth or Wi-Fi. These are battery-powered receiver systems, of which the headphone is only a component. Cordless headphones are used with events such as a Silent disco or Silent Gig.

In the professional audio sector, headphones are used in live situations by disc jockeys with a DJ mixer, and sound engineers for monitoring signal sources. In radio studios, DJs use a pair of headphones when talking to the microphone while the speakers are turned off to eliminate acoustic feedback while monitoring their own voice. In studio recordings, musicians and singers use headphones to play or sing along to a backing track or band. In military applications, audio signals of many varieties are monitored using headphones.

Wired headphones are attached to an audio source by a cable. The most common connectors are 6.35 mm (1/4 inch) and 3.5 mm phone connectors. The larger 6.35 mm connector is more common on fixed location home or professional equipment. The 3.5 mm connector remains the most widely used connector for portable application today. Adapters are available for converting between 6.35 mm and 3.5 mm devices.

As active component, wireless headphones tend to be costlier due to the necessity for internal hardware such as a battery, a charging controller, a speaker driver, and a wireless transceiver, whereas wired headphones are a passive component, outsourcing speaker driving to the audio source.

Some headphone cords are equipped with a serial potentiometer for volume control.

Wired headphones may be equipped with a non-detachable cable or a detachable auxiliary male-to-male plug, as well as some with two ports to allow connecting another wired headphone in a parallel circuit, which splits the audio signal to share with another participant, but can also be used to hear audio from two inputs simultaneously. An external audio splitter can retrofit this ability.

Applications for audiometric testing

Various types of specially designed headphones or earphones are also used to evaluate the status of the auditory system in the field of audiology for establishing hearing thresholds, medically diagnosing hearing loss, identifying other hearing related disease, and monitoring hearing status in occupational hearing conservation programs. Specific models of headphones have been adopted as the standard due to the ease of calibration and ability to compare results between testing facilities.

Supra-aural style headphones are historically the most commonly used in audiology as they are the easiest to calibrate and were considered the standard for many years. Commonly used models are the Telephonics Dynamic Headphone (TDH) 39, TDH-49, and TDH-50. In-the-ear or insert style earphones are used more commonly today as they provide higher levels of interaural attenuation, introduce less variability when testing 6,000 and 8,000 Hz, and avoid testing issues resulting from collapsed ear canals. A commonly used model of insert earphone is the Etymotic Research ER-3A. Circum-aural earphones are also used to establish hearing thresholds in the extended high frequency range (8,000 Hz to 20,000 kHz). Along with Etymotic Research ER-2A insert earphones, the Sennheiser HDA300 and Koss HV/1A circum-aural earphones are the only models that have reference equivalent threshold sound pressure level values for the extended high frequency range as described by ANSI standards.

Audiometers and headphones must be calibrated together. During the calibration process, the output signal from the audiometer to the headphones is measured with a sound level meter to ensure that the signal is accurate to the reading on the audiometer for sound pressure level and frequency. Calibration is done with the earphones in an acoustic coupler that is intended to mimic the transfer function of the outer ear. Because specific headphones are used in the initial audiometer calibration process, they cannot be replaced with any other set of headphones, even from the same make and model.

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Jai Ganesh · 2024-09-17 19:39:43

2306) Graphite

Gist

Graphite is a type of crystal carbon and a half-metal along with being one of the renowned carbon allotropes. Under the conditions that are ideal, it would be one of the most stable forms of carbon available. To define the standard state of heat for making compounds of carbons.

Graphite is used in pencils, lubricants, crucibles, foundry facings, polishes, brushes for electric motors, and cores of nuclear reactors.

Summary

Graphite is a mineral consisting of carbon. Graphite has a greasy feel and leaves a black mark, thus the name from the Greek verb graphein, “to write.”

Graphite has a layered structure that consists of rings of six carbon atoms arranged in widely spaced horizontal sheets. Graphite thus crystallizes in the hexagonal system, in contrast to diamond, another form of carbon, that crystallizes in the octahedral or tetrahedral system. Such pairs of differing forms of the same element usually are rather similar in their physical properties, but not so in this case. Graphite is dark gray to black, opaque, and very soft (with a Mohs scale hardness of 1.5), while diamond may be colorless and transparent and is the hardest naturally occurring substance (with a Mohs scale hardness of 10). Graphite is very soft because the individual layers of carbon atoms are not as tightly bound together as the atoms within the layer. It is an excellent conductor of heat and electricity.

Graphite is formed by the metamorphosis of sediments containing carbonaceous material, by the reaction of carbon compounds with hydrothermal solutions or magmatic fluids, or possibly by the crystallization of magmatic carbon. It occurs as isolated scales, large masses, or veins in older crystalline rocks, gneiss, schist, quartzite, and marble and also in granites, pegmatites, and carbonaceous clay slates. Small isometric crystals of graphitic carbon (possibly pseudomorphs after diamond) found in meteoritic iron are called cliftonite.

Naturally occurring graphite is classified into three types: amorphous, flake, and vein. Amorphous is the most common kind and is formed by metamorphism under low pressures and temperatures. It is found in coal and shale and has the lowest carbon content, typically 70 to 90 percent, of the three types. Flake graphite appears in flat layers and is formed by metamorphism under high pressures and temperatures. It is the most commonly used type and has a carbon content between 85 and 98 percent. Vein graphite is the rarest form and is likely formed when carbon compounds react with hydrothermal solutions or magmatic fluids. Vein graphite can have a purity greater than 99 percent and is commercially mined only in Sri Lanka.

Graphite was first synthesized accidentally by Edward G. Acheson while he was performing high-temperature experiments on carborundum. He found that at about 4,150 °C (7,500 °F) the silicon in the carborundum vaporized, leaving the carbon behind in graphitic form. Acheson was granted a patent for graphite manufacture in 1896, and commercial production started in 1897. Since 1918 petroleum coke, small and imperfect graphite crystals surrounded by organic compounds, has been the major raw material in the production of 99 to 99.5 percent pure graphite.

Graphite is used in pencils, lubricants, crucibles, foundry facings, polishes, brushes for electric motors, and cores of nuclear reactors. Its high thermal and electrical conductivity make it a key part of steelmaking, where it is used as electrodes in electric arc furnaces. In the early 21st century, global demand for graphite has increased because of its use as the anode in lithium-ion batteries for electric vehicles. About 75 percent of graphite is mined in China, with significant amounts mined in Madagascar, Mozambique, and Brazil.

Details

Graphite is a crystalline allotrope (form) of the element carbon. It consists of many stacked layers of graphene typically in the excess of hundred(s) of layers. Graphite occurs naturally and is the most stable form of carbon under standard conditions. Synthetic and natural graphite are consumed on a large scale (1.3 million metric tons per year in 2022) for uses in many critical industries including refractories (50%), lithium-ion batteries (18%), foundries (10%), lubricants (5%), among others (17%). Under extremely high pressures and extremely high temperatures it converts to diamond. It is a good conductor of both heat and electricity

Types and varieties:

Natural graphite

Graphite occurs naturally in ores that can be classified into one of two categories either amorphous (microcrystalline) or crystalline (flake or lump/chip) which is determined by the ore morphology, crystallinity, and grain size. All naturally occurring graphite deposits are formed from the metamorphism of carbonaceous sedimentary rocks, and the ore type is due to its geologic setting. Coal that has been thermally metamorphosed is the typical source of amorphous graphite. Crystalline flake graphite is mined from carbonaceous metamorphic rocks, while lump or chip graphite is mined from veins which occur in high-grade metamorphic regions. There are serious negative environmental impacts to graphite mining.

Synthetic graphite

Synthetic graphite is graphite of high purity produced by thermal graphitization at temperatures in excess of 2,100 °C from hydrocarbon materials most commonly by a process known as the Acheson process. The high temperatures are maintained for weeks, and are required not only to form the graphite from the precursor carbons but to also vaporize any impurities that may be present, including hydrogen, nitrogen, sulfur, organics, and metals. This is why synthetic graphite is highly pure in excess of 99.9% C purity, but typically has lower density, conductivity and a higher porosity than its natural equivalent. Synthetic graphite can also be formed into very large flakes (cm) while maintaining its high purity unlike almost all sources of natural graphite. Synthetic graphite has also been known to be formed by other methods including by chemical vapor deposition from hydrocarbons at temperatures above 2,500 K (2,230 °C), by decomposition of thermally unstable carbides or by crystallizing from metal melts supersaturated with carbon.

Biographite

Biographite is a commercial product proposal for reducing the carbon footprint of lithium iron phosphate (LFP) batteries. It is produced from forestry waste and similar byproducts by a company in New Zealand using a novel process called thermo-catalytic graphitisation which project is supported by grants from interested parties including a forestry company in Finland and a battery maker in Hong Kong.

Natural graphite:

Occurrence

Graphite occurs in metamorphic rocks as a result of the reduction of sedimentary carbon compounds during metamorphism. It also occurs in igneous rocks and in meteorites. Minerals associated with graphite include quartz, calcite, micas and tourmaline. The principal export sources of mined graphite are in order of tonnage: China, Mexico, Canada, Brazil, and Madagascar. Significant unexploited graphite resources also exists in Colombia's Cordillera Central in the form of graphite-bearing schists.

In meteorites, graphite occurs with troilite and silicate minerals. Small graphitic crystals in meteoritic iron are called cliftonite. Some microscopic grains have distinctive isotopic compositions, indicating that they were formed before the Solar System. They are one of about 12 known types of minerals that predate the Solar System and have also been detected in molecular clouds. These minerals were formed in the ejecta when supernovae exploded or low to intermediate-sized stars expelled their outer envelopes late in their lives. Graphite may be the second or third oldest mineral in the Universe.

Structure

Graphite consists of sheets of trigonal planar carbon. The individual layers are called graphene. In each layer, each carbon atom is bonded to three other atoms forming a continuous layer of sp2 bonded carbon hexagons, like a honeycomb lattice with a bond length of 0.142 nm, and the distance between planes is 0.335 nm. Bonding between layers is relatively weak van der Waals bonds, which allows the graphene-like layers to be easily separated and to glide past each other. Electrical conductivity perpendicular to the layers is consequently about 1000 times lower.

There are two allotropic forms called alpha (hexagonal) and beta (rhombohedral), differing in terms of the stacking of the graphene layers: stacking in alpha graphite is ABA, as opposed to ABC stacking in the energetically less stable beta graphite. Rhombohedral graphite cannot occur in pure form. Natural graphite, or commercial natural graphite, contains 5 to 15% rhombohedral graphite and this may be due to intensive milling. The alpha form can be converted to the beta form through shear forces, and the beta form reverts to the alpha form when it is heated to 1300 °C for four hours.

Thermodynamics

The equilibrium pressure and temperature conditions for a transition between graphite and diamond is well established theoretically and experimentally. The pressure changes linearly between 1.7 GPa at 0 K and 12 GPa at 5000 K (the diamond/graphite/liquid triple point). However, the phases have a wide region about this line where they can coexist. At normal temperature and pressure, 20 °C (293 K) and 1 standard atmosphere (0.10 MPa), the stable phase of carbon is graphite, but diamond is metastable and its rate of conversion to graphite is negligible. However, at temperatures above about 4500 K, diamond rapidly converts to graphite. Rapid conversion of graphite to diamond requires pressures well above the equilibrium line: at 2000 K, a pressure of 35 GPa is needed.

Other properties

The acoustic and thermal properties of graphite are highly anisotropic, since phonons propagate quickly along the tightly bound planes, but are slower to travel from one plane to another. Graphite's high thermal stability and electrical and thermal conductivity facilitate its widespread use as electrodes and refractories in high temperature material processing applications. However, in oxygen-containing atmospheres graphite readily oxidizes to form carbon dioxide at temperatures of 700 °C and above.

Graphite is an electrical conductor, hence useful in such applications as arc lamp electrodes. It can conduct electricity due to the vast electron delocalization within the carbon layers (a phenomenon called aromaticity). These valence electrons are free to move, so are able to conduct electricity. However, the electricity is primarily conducted within the plane of the layers. The conductive properties of powdered graphite allow its use as pressure sensor in carbon microphones.

Graphite and graphite powder are valued in industrial applications for their self-lubricating and dry lubricating properties. However, the use of graphite is limited by its tendency to facilitate pitting corrosion in some stainless steel, and to promote galvanic corrosion between dissimilar metals (due to its electrical conductivity). It is also corrosive to aluminium in the presence of moisture. For this reason, the US Air Force banned its use as a lubricant in aluminium aircraft, and discouraged its use in aluminium-containing automatic weapons. Even graphite pencil marks on aluminium parts may facilitate corrosion. Another high-temperature lubricant, hexagonal boron nitride, has the same molecular structure as graphite. It is sometimes called white graphite, due to its similar properties.

When a large number of crystallographic defects bind these planes together, graphite loses its lubrication properties and becomes what is known as pyrolytic graphite. It is also highly anisotropic, and diamagnetic, thus it will float in mid-air above a strong magnet. (If it is made in a fluidized bed at 1000–1300 °C then it is isotropic turbostratic, and is used in blood-contacting devices like mechanical heart valves and is called pyrolytic carbon, and is not diamagnetic. Pyrolytic graphite and pyrolytic carbon are often confused but are very different materials.)

Natural and crystalline graphites are not often used in pure form as structural materials, due to their shear-planes, brittleness, and inconsistent mechanical properties.

Additional Information

Graphite is a mineral composed of stacked sheets of carbon atoms with a hexagonal crystal structure. It is the most stable form of pure carbon under standard conditions. Graphite is very soft, has a low specific gravity, is relatively non-reactive, and has high electrical and thermal conductivity.

Graphite occurs naturally in igneous and metamorphic rocks, where high temperatures and pressures compress carbon into graphite. Graphite can also be created synthetically by heating materials with high carbon content (e.g. petroleum coke or coal-tar pitch). The carbon-rich material is heated to 2500 to 3000 degrees Celsius, which is hot enough to "purify" the material of contaminants, allowing the carbon to form its hexagonal sheets.

Graphite is extremely soft and breaks into thin flexible flakes that easily slide over one another, resulting in a greasy feel. Due to this, graphite is a good "dry" lubricant and can be used in applications where wet lubricants (like lubricating oil) cannot.

Carbon has several other allotropes, or forms, that occur naturally, each with their own crystal structure. One form is graphene, which is a single layer of carbon atoms in a hexagonal pattern. Another well-known allotrope of carbon, are diamonds. Although also composed of pure carbon, diamonds are almost entirely different in their physical properties.

Uses

Graphite is used in a number of applications that require high temperatures and need a material that will not melt or disintegrate. Graphite is used to make the crucibles for the steel industry. Graphite is also used as a neutron moderator in certain nuclear reactors, like the Soviet RBMK, due to its ability to slow down fast-moving neutrons.

Other common uses of graphite include:

* Pencil lead
* Lubricant
* Electrodes in batteries
* Brake linings for heavy vehicles.

Jai Ganesh · Yesterday 00:18:50

2307) Ulna

Gist

The ulna is one of the two forearm long bones that, in conjunction with the radius, make up the antebrachium. The bone spans from the elbow to the wrist on the medial side of the forearm when in anatomical position. In comparison to the radius, the ulna is described to be larger and longer.

Summary

Ulna is inner of two bones of the forearm when viewed with the palm facing forward. (The other, shorter bone of the forearm is the radius.) The upper end of the ulna presents a large C-shaped notch—the semilunar, or trochlear, notch—which articulates with the trochlea of the humerus (upper arm bone) to form the elbow joint. The projection that forms the upper border of this notch is called the olecranon process; it articulates behind the humerus in the olecranon fossa and may be felt as the point of the elbow. The projection that forms the lower border of the trochlear notch, the coronoid process, enters the coronoid fossa of the humerus when the elbow is flexed. On the outer side is the radial notch, which articulates with the head of the radius. The head of the bone is elsewhere roughened for muscle attachment. The shaft is triangular in cross section; an interosseous ridge extends its length and provides attachment for the interosseous membrane connecting the ulna and the radius. The lower end of the bone presents a small cylindrical head that articulates with the radius at the side and the wrist bones below. Also at the lower end is a styloid process, medially, that articulates with a disk between it and the cuneiform (os triquetrum) wrist bone.

The ulna is present in all land vertebrates. In amphibians and some reptiles the radius and ulna do not articulate. The elbow joint evolved first among birds and mammals. The radius tends to be slender in birds; but the ulna is more often reduced in mammals, especially in those adapted for running and, in the case of bats, flying.

Details:

The ulna or ulnar bone (pl.: ulnae or ulnas) is a long bone in the forearm stretching from the elbow to the wrist. It is on the same side of the forearm as the little finger, running parallel to the radius, the forearm's other long bone. Longer and thinner than the radius, the ulna is considered to be the smaller long bone of the lower arm. The corresponding bone in the lower leg is the fibula.

Structure

The ulna is a long bone found in the forearm that stretches from the elbow to the wrist, and when in standard anatomical position, is found on the medial side of the forearm. It is broader close to the elbow, and narrows as it approaches the wrist.

Close to the elbow, the ulna has a bony process, the olecranon process, a hook-like structure that fits into the olecranon fossa of the humerus. This prevents hyperextension and forms a hinge joint with the trochlea of the humerus. There is also a radial notch for the head of the radius, and the ulnar tuberosity to which muscles attach.

Close to the wrist, the ulna has a styloid process.

Near the elbow

Near the elbow, the ulna has two curved processes, the olecranon and the coronoid process; and two concave, articular cavities, the semilunar and radial notches.

The olecranon is a large, thick, curved eminence, situated at the upper and back part of the ulna. It is bent forward at the summit so as to present a prominent lip which is received into the olecranon fossa of the humerus in extension of the forearm. Its base is contracted where it joins the body and the narrowest part of the upper end of the ulna. Its posterior surface, directed backward, is triangular, smooth, subcutaneous, and covered by a bursa. Its superior surface is of quadrilateral form, marked behind by a rough impression for the insertion of the triceps brachii; and in front, near the margin, by a slight transverse groove for the attachment of part of the posterior ligament of the elbow joint. Its anterior surface is smooth, concave, and forms the upper part of the semilunar notch. Its borders present continuations of the groove on the margin of the superior surface; they serve for the attachment of ligaments: the back part of the ulnar collateral ligament medially, and the posterior ligament laterally. From the medial border a part of the flexor carpi ulnaris arises; while to the lateral border the anconeus is attached.

The coronoid process is a triangular eminence projecting forward from the upper and front part of the ulna. Its base is continuous with the body of the bone, and of considerable strength. Its apex is pointed, slightly curved upward, and in flexion of the forearm is received into the coronoid fossa of the humerus. Its upper surface is smooth, concave, and forms the lower part of the semilunar notch. Its antero-inferior surface is concave, and marked by a rough impression for the insertion of the brachialis. At the junction of this surface with the front of the body is a rough eminence, the tuberosity of the ulna, which gives insertion to a part of the brachialis; to the lateral border of this tuberosity the oblique cord is attached. Its lateral surface presents a narrow, oblong, articular depression, the radial notch. Its medial surface, by its prominent, free margin, serves for the attachment of part of the ulnar collateral ligament. At the front part of this surface is a small rounded eminence for the origin of one head of the flexor digitorum superficialis; behind the eminence is a depression for part of the origin of the flexor digitorum profundus; descending from the eminence is a ridge which gives origin to one head of the pronator teres. Frequently, the flexor pollicis longus arises from the lower part of the coronoid process by a rounded bundle of muscular fibers.

The semilunar notch is a large depression, formed by the olecranon and the coronoid process, and serving as articulation with the trochlea of the humerus. About the middle of either side of this notch is an indentation, which contracts it somewhat, and indicates the junction of the olecranon and the coronoid process. The notch is concave from above downward, and divided into a medial and a lateral portion by a smooth ridge running from the summit of the olecranon to the tip of the coronoid process. The medial portion is the larger, and is slightly concave transversely; the lateral is convex above, slightly concave below.

The radial notch is a narrow, oblong, articular depression on the lateral side of the coronoid process; it receives the circumferential articular surface of the head of the radius. It is concave from before backward, and its prominent extremities serve for the attachment of the annular ligament.

Body

The body of the ulna at its upper part is prismatic in form, and curved so as to be convex behind and lateralward; its central part is straight; its lower part is rounded, smooth, and bent a little lateralward. It tapers gradually from above downward, and has three borders and three surfaces.

Borders

* The volar border (margo volaris; anterior border) begins above at the prominent medial angle of the coronoid process, and ends below in front of the styloid process. Its upper part, well-defined, and its middle portion, smooth and rounded, give origin to the flexor digitorum profundus; its lower fourth serves for the origin of the pronator quadratus. This border separates the volar from the medial surface.
* The dorsal border (margo dorsalis; posterior border) begins above at the apex of the triangular subcutaneous surface at the back part of the olecranon, and ends below at the back of the styloid process; it is well-marked in the upper three-fourths, and gives attachment to an aponeurosis which affords a common origin to the flexor carpi ulnaris, the extensor carpi ulnaris, and the flexor digitorum profundus; its lower fourth is smooth and rounded. This border separates the medial from the dorsal surface.
* The interosseous crest (crista interossea; external or interosseous border) begins above by the union of two lines, which converge from the extremities of the radial notch and enclose between them a triangular space for the origin of part of the supinator; it ends below at the head of the ulna. Its upper part is sharp, its lower fourth smooth and rounded. This crest gives attachment to the interosseous membrane, and separates the volar from the dorsal surface.

Surfaces

* The volar surface (facies volaris; anterior surface), much broader above than below, is concave in its upper three-fourths, and gives origin to the flexor digitorum profundus; its lower fourth, also concave, is covered by the pronator quadratus. The lower fourth is separated from the remaining portion by a ridge, directed obliquely downward and medialward, which marks the extent of origin of the pronator quadratus. At the junction of the upper with the middle third of the bone is the nutrient canal, directed obliquely upward.
* The dorsal surface (facies dorsalis; posterior surface) directed backward and lateralward, is broad and concave above; convex and somewhat narrower in the middle; narrow, smooth, and rounded below. On its upper part is an oblique ridge, which runs from the dorsal end of the radial notch, downward to the dorsal border; the triangular surface above this ridge receives the insertion of the anconeus, while the upper part of the ridge affords attachment to the supinator. Below this the surface is subdivided by a longitudinal ridge, sometimes called the perpendicular line, into two parts: the medial part is smooth, and covered by the extensor carpi ulnaris; the lateral portion, wider and rougher, gives origin from above downward to the supinator, the abductor pollicis longus, the extensor pollicis longus, and the extensor indicis proprius.
* The medial surface (facies medialis; internal surface) is broad and concave above, narrow and convex below. Its upper three-fourths give origin to the flexor digitorum profundus; its lower fourth is subcutaneous.

Near the wrist

Near the wrist, the ulnar, with two eminences; the lateral and larger is a rounded, articular eminence, termed the head of the ulna; the medial, narrower and more projecting, is a non-articular eminence, the ulnar styloid process.

* The head of the ulna presents an articular surface, part of which, of an oval or semilunar form, is directed downward, and articulates with the upper surface of the triangular articular disk which separates it from the wrist-joint; the remaining portion, directed lateralward, is narrow, convex, and received into the ulnar notch of the radius.
* The styloid process projects from the medial and back part of the bone; it descends a little lower than the head, and its rounded end affords attachment to the ulnar collateral ligament of the wrist-joint.
* The head is separated from the styloid process by a depression for the attachment of the apex of the triangular articular disk, and behind, by a shallow groove for the tendon of the extensor carpi ulnaris.

Microanatomy

The ulna is a long bone. The long, narrow medullary cavity of the ulna is enclosed in a strong wall of cortical tissue which is thickest along the interosseous border and dorsal surface. At the extremities the compact layer thins. The compact layer is continued onto the back of the olecranon as a plate of close spongy bone with lamellae parallel. From the inner surface of this plate and the compact layer below it trabeculae arch forward toward the olecranon and coronoid and cross other trabeculae, passing backward over the medullary cavity from the upper part of the shaft below the coronoid. Below the coronoid process there is a small area of compact bone from which trabeculae curve upward to end obliquely to the surface of the semilunar notch which is coated with a thin layer of compact bone. The trabeculae at the lower end have a more longitudinal direction.

Development

The ulna is ossified from three centers: one each for the body, the wrist end, and the elbow end, near the top of the olecranon. Ossification begins near the middle of the body of the ulna, about the eighth week of fetal life, and soon extends through the greater part of the bone.

At birth, the ends are cartilaginous. About the fourth year or so, a center appears in the middle of the head, and soon extends into the ulnar styloid process. About the tenth year, a center appears in the olecranon near its extremity, the chief part of this process being formed by an upward extension of the body. The upper epiphysis joins the body about the sixteenth, the lower about the twentieth year.

Function:

Joints

The ulna forms part of the wrist joint and elbow joints. Specifically, the ulna joins (articulates) with:

* trochlea of the humerus, at the right side elbow as a hinge joint with semilunar trochlear notch of the ulna.
* the radius, near the elbow as a pivot joint, this allows the radius to cross over the ulna in pronation.
* the distal radius, where it fits into the ulnar notch.
* the radius along its length via the interosseous membrane that forms a syndesmosis joint.

Additional Information

The ulna is a long thin bone with a small distal head that bears the styloid process, and an expanded proximal end. The proximal end terminates in the olecranon process and bears the semilunar notch on its upper surface. In man, the head of the ulna does not articulate with any of the bones of the carpus. In the rat, the ulna may articulate via its styloid process with the triquetrum and possibly the pisiform bone (see below for a discussion of this point). The radius and ulna are connected more or less throughout their length by an interosseus ligament, which contributes to the origins of some muscles of the forearm.

Jai Ganesh · Yesterday 21:50:31

2308) Battery

Gist

A battery is a device that converts chemical energy contained within its active materials directly into electric energy by means of an electrochemical oxidation-reduction (redox) reaction. This type of reaction involves the transfer of electrons from one material to another via an electric circuit.

Batteries consist of two electrical terminals called the cathode and the anode, separated by a chemical material called an electrolyte. To accept and release energy, a battery is coupled to an external circuit.

A chemical reaction between the metals and the electrolyte frees more electrons in one metal than it does in the other. The metal that frees more electrons develops a positive charge, and the other metal develops a negative charge.

Summary

A battery, in electricity and electrochemistry, is any of a class of devices that convert chemical energy directly into electrical energy. Although the term battery, in strict usage, designates an assembly of two or more galvanic cells capable of such energy conversion, it is commonly applied to a single cell of this kind.

Every battery (or cell) has a cathode, or positive plate, and an anode, or negative plate. These electrodes must be separated by and are often immersed in an electrolyte that permits the passage of ions between the electrodes. The electrode materials and the electrolyte are chosen and arranged so that sufficient electromotive force (measured in volts) and electric current (measured in amperes) can be developed between the terminals of a battery to operate lights, machines, or other devices. Since an electrode contains only a limited number of units of chemical energy convertible to electrical energy, it follows that a battery of a given size has only a certain capacity to operate devices and will eventually become exhausted. The active parts of a battery are usually encased in a box with a cover system (or jacket) that keeps air outside and the electrolyte solvent inside and that provides a structure for the assembly.

Commercially available batteries are designed and built with market factors in mind. The quality of materials and the complexity of electrode and container design are reflected in the market price sought for any specific product. As new materials are discovered or the properties of traditional ones improved, however, the typical performance of even older battery systems sometimes increases by large percentages.

Batteries are divided into two general groups: (1) primary batteries and (2) secondary, or storage, batteries. Primary batteries are designed to be used until the voltage is too low to operate a given device and are then discarded. Secondary batteries have many special design features, as well as particular materials for the electrodes, that permit them to be reconstituted (recharged). After partial or complete discharge, they can be recharged by the application of direct current (DC) voltage. While the original state is usually not restored completely, the loss per recharging cycle in commercial batteries is only a small fraction of 1 percent even under varied conditions.

Details

An electric battery is a source of electric power consisting of one or more electrochemical cells with external connections for powering electrical devices. When a battery is supplying power, its positive terminal is the cathode and its negative terminal is the anode. The terminal marked negative is the source of electrons that will flow through an external electric circuit to the positive terminal. When a battery is connected to an external electric load, a redox reaction converts high-energy reactants to lower-energy products, and the free-energy difference is delivered to the external circuit as electrical energy. Historically the term "battery" specifically referred to a device composed of multiple cells; however, the usage has evolved to include devices composed of a single cell.

Primary (single-use or "disposable") batteries are used once and discarded, as the electrode materials are irreversibly changed during discharge; a common example is the alkaline battery used for flashlights and a multitude of portable electronic devices. Secondary (rechargeable) batteries can be discharged and recharged multiple times using an applied electric current; the original composition of the electrodes can be restored by reverse current. Examples include the lead–acid batteries used in vehicles and lithium-ion batteries used for portable electronics such as laptops and mobile phones.

Batteries come in many shapes and sizes, from miniature cells used to power hearing aids and wristwatches to, at the largest extreme, huge battery banks the size of rooms that provide standby or emergency power for telephone exchanges and computer data centers. Batteries have much lower specific energy (energy per unit mass) than common fuels such as gasoline. In automobiles, this is somewhat offset by the higher efficiency of electric motors in converting electrical energy to mechanical work, compared to combustion engines.

History:

Invention

Benjamin Franklin first used the term "battery" in 1749 when he was doing experiments with electricity using a set of linked Leyden jar capacitors. Franklin grouped a number of the jars into what he described as a "battery", using the military term for weapons functioning together. By multiplying the number of holding vessels, a stronger charge could be stored, and more power would be available on discharge.

Italian physicist Alessandro Volta built and described the first electrochemical battery, the voltaic pile, in 1800. This was a stack of copper and zinc plates, separated by brine-soaked paper disks, that could produce a steady current for a considerable length of time. Volta did not understand that the voltage was due to chemical reactions. He thought that his cells were an inexhaustible source of energy, and that the associated corrosion effects at the electrodes were a mere nuisance, rather than an unavoidable consequence of their operation, as Michael Faraday showed in 1834.

Although early batteries were of great value for experimental purposes, in practice their voltages fluctuated and they could not provide a large current for a sustained period. The Daniell cell, invented in 1836 by British chemist John Frederic Daniell, was the first practical source of electricity, becoming an industry standard and seeing widespread adoption as a power source for electrical telegraph networks. It consisted of a copper pot filled with a copper sulfate solution, in which was immersed an unglazed earthenware container filled with sulfuric acid and a zinc electrode.

These wet cells used liquid electrolytes, which were prone to leakage and spillage if not handled correctly. Many used glass jars to hold their components, which made them fragile and potentially dangerous. These characteristics made wet cells unsuitable for portable appliances. Near the end of the nineteenth century, the invention of dry cell batteries, which replaced the liquid electrolyte with a paste, made portable electrical devices practical.

Batteries in vacuum tube devices historically used a wet cell for the "A" battery (to provide power to the filament) and a dry cell for the "B" battery (to provide the plate voltage).

Future

Between 2010 and 2018, annual battery demand grew by 30%, reaching a total of 180 GWh in 2018. Conservatively, the growth rate is expected to be maintained at an estimated 25%, culminating in demand reaching 2600 GWh in 2030. In addition, cost reductions are expected to further increase the demand to as much as 3562 GWh.

Important reasons for this high rate of growth of the electric battery industry include the electrification of transport, and large-scale deployment in electricity grids, supported by anthropogenic climate change-driven moves away from fossil-fuel combusted energy sources to cleaner, renewable sources, and more stringent emission regimes.

Distributed electric batteries, such as those used in battery electric vehicles (vehicle-to-grid), and in home energy storage, with smart metering and that are connected to smart grids for demand response, are active participants in smart power supply grids. New methods of reuse, such as echelon use of partly-used batteries, add to the overall utility of electric batteries, reduce energy storage costs, and also reduce pollution/emission impacts due to longer lives. In echelon use of batteries, vehicle electric batteries that have their battery capacity reduced to less than 80%, usually after service of 5–8 years, are repurposed for use as backup supply or for renewable energy storage systems.

Grid scale energy storage envisages the large-scale use of batteries to collect and store energy from the grid or a power plant and then discharge that energy at a later time to provide electricity or other grid services when needed. Grid scale energy storage (either turnkey or distributed) are important components of smart power supply grids.

Chemistry and principles

Batteries convert chemical energy directly to electrical energy. In many cases, the electrical energy released is the difference in the cohesive or bond energies of the metals, oxides, or molecules undergoing the electrochemical reaction. For instance, energy can be stored in Zn or Li, which are high-energy metals because they are not stabilized by d-electron bonding, unlike transition metals. Batteries are designed so that the energetically favorable redox reaction can occur only when electrons move through the external part of the circuit.

A battery consists of some number of voltaic cells. Each cell consists of two half-cells connected in series by a conductive electrolyte containing metal cations. One half-cell includes electrolyte and the negative electrode, the electrode to which anions (negatively charged ions) migrate; the other half-cell includes electrolyte and the positive electrode, to which cations (positively charged ions) migrate. Cations are reduced (electrons are added) at the cathode, while metal atoms are oxidized (electrons are removed) at the anode. Some cells use different electrolytes for each half-cell; then a separator is used to prevent mixing of the electrolytes while allowing ions to flow between half-cells to complete the electrical circuit.

The voltage developed across a cell's terminals depends on the energy release of the chemical reactions of its electrodes and electrolyte. Alkaline and zinc–carbon cells have different chemistries, but approximately the same emf of 1.5 volts; likewise NiCd and NiMH cells have different chemistries, but approximately the same emf of 1.2 volts. The high electrochemical potential changes in the reactions of lithium compounds give lithium cells emfs of 3 volts or more.

Almost any liquid or moist object that has enough ions to be electrically conductive can serve as the electrolyte for a cell. As a novelty or science demonstration, it is possible to insert two electrodes made of different metals into a lemon, potato, etc. and generate small amounts of electricity.

A voltaic pile can be made from two coins (such as a nickel and a penny) and a piece of paper towel dipped in salt water. Such a pile generates a very low voltage but, when many are stacked in series, they can replace normal batteries for a short time.

Types:

Primary and secondary batteries

Batteries are classified into primary and secondary forms:

* Primary batteries are designed to be used until exhausted of energy then discarded. Their chemical reactions are generally not reversible, so they cannot be recharged. When the supply of reactants in the battery is exhausted, the battery stops producing current and is useless.
* Secondary batteries can be recharged; that is, they can have their chemical reactions reversed by applying electric current to the cell. This regenerates the original chemical reactants, so they can be used, recharged, and used again multiple times.

Some types of primary batteries used, for example, for telegraph circuits, were restored to operation by replacing the electrodes. Secondary batteries are not indefinitely rechargeable due to dissipation of the active materials, loss of electrolyte and internal corrosion.

Primary batteries, or primary cells, can produce current immediately on assembly. These are most commonly used in portable devices that have low current drain, are used only intermittently, or are used well away from an alternative power source, such as in alarm and communication circuits where other electric power is only intermittently available. Disposable primary cells cannot be reliably recharged, since the chemical reactions are not easily reversible and active materials may not return to their original forms. Battery manufacturers recommend against attempting to recharge primary cells. In general, these have higher energy densities than rechargeable batteries, but disposable batteries do not fare well under high-drain applications with loads under 75 ohms. Common types of disposable batteries include zinc–carbon batteries and alkaline batteries.

Secondary batteries, also known as secondary cells, or rechargeable batteries, must be charged before first use; they are usually assembled with active materials in the discharged state. Rechargeable batteries are (re)charged by applying electric current, which reverses the chemical reactions that occur during discharge/use. Devices to supply the appropriate current are called chargers. The oldest form of rechargeable battery is the lead–acid battery, which are widely used in automotive and boating applications. This technology contains liquid electrolyte in an unsealed container, requiring that the battery be kept upright and the area be well ventilated to ensure safe dispersal of the hydrogen gas it produces during overcharging. The lead–acid battery is relatively heavy for the amount of electrical energy it can supply. Its low manufacturing cost and its high surge current levels make it common where its capacity (over approximately 10 Ah) is more important than weight and handling issues. A common application is the modern car battery, which can, in general, deliver a peak current of 450 amperes.

Composition

Many types of electrochemical cells have been produced, with varying chemical processes and designs, including galvanic cells, electrolytic cells, fuel cells, flow cells and voltaic piles.

A wet cell battery has a liquid electrolyte. Other names are flooded cell, since the liquid covers all internal parts or vented cell, since gases produced during operation can escape to the air. Wet cells were a precursor to dry cells and are commonly used as a learning tool for electrochemistry. They can be built with common laboratory supplies, such as beakers, for demonstrations of how electrochemical cells work. A particular type of wet cell known as a concentration cell is important in understanding corrosion. Wet cells may be primary cells (non-rechargeable) or secondary cells (rechargeable). Originally, all practical primary batteries such as the Daniell cell were built as open-top glass jar wet cells. Other primary wet cells are the Leclanche cell, Grove cell, Bunsen cell, Chromic acid cell, Clark cell, and Weston cell. The Leclanche cell chemistry was adapted to the first dry cells. Wet cells are still used in automobile batteries and in industry for standby power for switchgear, telecommunication or large uninterruptible power supplies, but in many places batteries with gel cells have been used instead. These applications commonly use lead–acid or nickel–cadmium cells. Molten salt batteries are primary or secondary batteries that use a molten salt as electrolyte. They operate at high temperatures and must be well insulated to retain heat.

A dry cell uses a paste electrolyte, with only enough moisture to allow current to flow. Unlike a wet cell, a dry cell can operate in any orientation without spilling, as it contains no free liquid, making it suitable for portable equipment. By comparison, the first wet cells were typically fragile glass containers with lead rods hanging from the open top and needed careful handling to avoid spillage. Lead–acid batteries did not achieve the safety and portability of the dry cell until the development of the gel battery. A common dry cell is the zinc–carbon battery, sometimes called the dry Leclanché cell, with a nominal voltage of 1.5 volts, the same as the alkaline battery (since both use the same zinc–manganese dioxide combination). A standard dry cell comprises a zinc anode, usually in the form of a cylindrical pot, with a carbon cathode in the form of a central rod. The electrolyte is ammonium chloride in the form of a paste next to the zinc anode. The remaining space between the electrolyte and carbon cathode is taken up by a second paste consisting of ammonium chloride and manganese dioxide, the latter acting as a depolariser. In some designs, the ammonium chloride is replaced by zinc chloride.

A reserve battery can be stored unassembled (unactivated and supplying no power) for a long period (perhaps years). When the battery is needed, then it is assembled (e.g., by adding electrolyte); once assembled, the battery is charged and ready to work. For example, a battery for an electronic artillery fuze might be activated by the impact of firing a gun. The acceleration breaks a capsule of electrolyte that activates the battery and powers the fuze's circuits. Reserve batteries are usually designed for a short service life (seconds or minutes) after long storage (years). A water-activated battery for oceanographic instruments or military applications becomes activated on immersion in water.

On 28 February 2017, the University of Texas at Austin issued a press release about a new type of solid-state battery, developed by a team led by lithium-ion battery inventor John Goodenough, "that could lead to safer, faster-charging, longer-lasting rechargeable batteries for handheld mobile devices, electric cars and stationary energy storage". The solid-state battery is also said to have "three times the energy density", increasing its useful life in electric vehicles, for example. It should also be more ecologically sound since the technology uses less expensive, earth-friendly materials such as sodium extracted from seawater. They also have much longer life.

Sony has developed a biological battery that generates electricity from sugar in a way that is similar to the processes observed in living organisms. The battery generates electricity through the use of enzymes that break down carbohydrates.

The sealed valve regulated lead–acid battery (VRLA battery) is popular in the automotive industry as a replacement for the lead–acid wet cell. The VRLA battery uses an immobilized sulfuric acid electrolyte, reducing the chance of leakage and extending shelf life. VRLA batteries immobilize the electrolyte. The two types are:

* Gel batteries (or "gel cell") use a semi-solid electrolyte.
* Absorbed Glass Mat (AGM) batteries absorb the electrolyte in a special fiberglass matting.

Other portable rechargeable batteries include several sealed "dry cell" types, that are useful in applications such as mobile phones and laptop computers. Cells of this type (in order of increasing power density and cost) include nickel–cadmium (NiCd), nickel–zinc (NiZn), nickel–metal hydride (NiMH), and lithium-ion (Li-ion) cells. Li-ion has by far the highest share of the dry cell rechargeable market. NiMH has replaced NiCd in most applications due to its higher capacity, but NiCd remains in use in power tools, two-way radios, and medical equipment.

In the 2000s, developments include batteries with embedded electronics such as USBCELL, which allows charging an AA battery through a USB connector, nanoball batteries that allow for a discharge rate about 100x greater than current batteries, and smart battery packs with state-of-charge monitors and battery protection circuits that prevent damage on over-discharge. Low self-discharge (LSD) allows secondary cells to be charged prior to shipping.

Lithium–sulfur batteries were used on the longest and highest solar-powered flight.

Consumer and industrial grades

Batteries of all types are manufactured in consumer and industrial grades. Costlier industrial-grade batteries may use chemistries that provide higher power-to-size ratio, have lower self-discharge and hence longer life when not in use, more resistance to leakage and, for example, ability to handle the high temperature and humidity associated with medical autoclave sterilization.

Additional Information:

Battery

A battery is a device that stores energy and then discharges it by converting chemical energy into electricity. Typical batteries most often produce electricity by chemical means through the use of one or more electrochemical cells. Many different materials can and have been used in batteries, but the common battery types are alkaline, lithium-ion, lithium-polymer, and nickel-metal hydride. Batteries can be connected to each other in a series circuit or a parallel circuit.

There is a wide variety of batteries that are available for purchase, and these different types of batteries are used in different devices. Large batteries are used to start cars, while much smaller batteries can power hearing aids. Overall, batteries are extremely important in everyday life.

Cells

A cell is a single unit that produces electricity through some method. Generally speaking, cells generate power through a thermal, chemical or optical process.

A typical cell has two terminals (referred to as electrodes) immersed in a chemical (referred to as the electrolyte). The two electrodes are separated by a porous wall or bridge which allows electric charge to pass from one side to the other through the electrolyte. The anode—the negative terminal—gains electrons while the cathode—the positive terminal—loses electrons. This exchange of electrons allows a difference in potential or voltage difference to be developed between the two terminals—allowing electricity to flow.

There can be a vast number of cells in a battery, from a single cell in an AA battery, to more than 7,100 cells in the 85 kWh Tesla Model S battery.

Primary cells ("dry")

In these cells a chemical action between the electrodes and electrolyte causes a permanent change, meaning they are not rechargeable. These batteries are single use, which results in more waste from the use of these batteries since they are disposed of after a relatively short period of time.

Secondary cells ("wet")

This type of cell (referred to as wet due to using a liquid electrolyte) generates a current through a secondary cell in the opposite direction of the first/normal cell. This causes the chemical action to go in reverse, effectively being restored, meaning that they are rechargeable. These batteries can be more expensive to purchase but generate less waste as they can be used several times.

Battery Capacity

Batteries are often rated in terms of their output voltage and capacity. The capacity is how long a particular battery will last in Ah (Ampere hours):

* A battery with a capacity of 1 Ah will last for one hour operating at 1 A.

Batteries can also be rated by their energy capacity. This is either done in watt-hours or kilowatt-hours.

* A battery with a capacity of 1 kWh will last for one hour while outputting 1 kW of electricity.

Jai Ganesh · Today 00:02:34

2309) Precipitation (chemistry)

Gist

In an aqueous solution, precipitation is the "sedimentation of a solid material (a precipitate) from a liquid solution". The solid formed is called the precipitate. In case of an inorganic chemical reaction leading to precipitation, the chemical reagent causing the solid to form is called the precipitant.

Summary

Chemical precipitation is the most common technology used in removing dissolved (ionic) metals from solutions, such as process wastewaters containing toxic metals. The ionic metals are converted to an insoluble form (particle) by the chemical reaction between the soluble metal compounds and the precipitating reagent. The particles formed by this reaction are removed from solution by settling and/or filtration. The unit operations typically required in this technology include neutralization, precipitation, coagulation/flocculation, solids/liquid separation, and dewatering.

The effectiveness of a chemical precipitation process is dependent on several factors, including the type and concentration of ionic metals present in solution, the precipitant used, the reaction conditions (especially the pH of the solution), and the presence of other constituents that may inhibit the precipitation reaction.

The most widely used chemical precipitation process is hydroxide precipitation (also referred to as precipitation by pH), in which metal hydroxides are formed by using calcium hydroxide (lime) or sodium hydroxide (caustic) as the precipitant. Each dissolved metal has a distinct pH value at which the optimum hydroxide precipitation occurs—from 7.5 for chromium to 11.0 for cadmium. Metal hydroxides are amphoteric, which means that they are increasingly soluble at both low and high pH values. Therefore, the optimum pH for precipitation of one metal may cause another metal to solubilize or start to go back into solution.

Details

In an aqueous solution, precipitation is the "sedimentation of a solid material (a precipitate) from a liquid solution". The solid formed is called the precipitate. In case of an inorganic chemical reaction leading to precipitation, the chemical reagent causing the solid to form is called the precipitant.

The clear liquid remaining above the precipitated or the centrifuged solid phase is also called the supernate or supernatant.

The notion of precipitation can also be extended to other domains of chemistry (organic chemistry and biochemistry) and even be applied to the solid phases (e.g. metallurgy and alloys) when solid impurities segregate from a solid phase.

Supersaturation

The precipitation of a compound may occur when its concentration exceeds its solubility. This can be due to temperature changes, solvent evaporation, or by mixing solvents. Precipitation occurs more rapidly from a strongly supersaturated solution.

The formation of a precipitate can be caused by a chemical reaction. When a barium chloride solution reacts with sulphuric acid, a white precipitate of barium sulfate is formed. When a potassium iodide solution reacts with a lead(II) nitrate solution, a yellow precipitate of lead(II) iodide is formed.

Inorganic chemistry

Precipitate formation is useful in the detection of the type of cation in a salt. To do this, an alkali first reacts with the unknown salt to produce a precipitate that is the hydroxide of the unknown salt. To identify the cation, the color of the precipitate and its solubility in excess are noted. Similar processes are often used in sequence – for example, a barium nitrate solution will react with sulfate ions to form a solid barium sulfate precipitate, indicating that it is likely that sulfate ions are present.

A common example of precipitation from aqueous solution is that of silver chloride. When silver nitrate (AgNO3) is added to a solution of potassium chloride (KCl) the precipitation of a white solid (AgCl) is observed.

Colloidal suspensions

Without sufficient attraction forces (e.g., Van der Waals force) to aggregate the solid particles together and to remove them from solution by gravity (settling), they remain in suspension and form colloids. Sedimentation can be accelerated by high speed centrifugation. The compact mass thus obtained is sometimes referred to as a 'pellet'.

Digestion and precipitates ageing

Digestion, or precipitate ageing, happens when a freshly formed precipitate is left, usually at a higher temperature, in the solution from which it precipitates. It results in purer and larger recrystallized particles. The physico-chemical process underlying digestion is called Ostwald ripening.

Organic chemistry

While precipitation reactions can be used for making pigments, removing ions from solution in water treatment, and in classical qualitative inorganic analysis, precipitation is also commonly used to isolate the products of an organic reaction during workup and purification operations. Ideally, the product of the reaction is insoluble in the solvent used for the reaction. Thus, it precipitates as it is formed, preferably forming pure crystals. An example of this would be the synthesis of porphyrins in refluxing propionic acid. By cooling the reaction mixture to room temperature, crystals of the porphyrin precipitate, and are collected by filtration on a Büchner filter.

Porphyrin synthesis

Precipitation may also occur when an antisolvent (a solvent in which the product is insoluble) is added, drastically reducing the solubility of the desired product. Thereafter, the precipitate may be easily separated by decanting, filtration, or by centrifugation. An example would be the synthesis of Cr3+tetraphenylporphyrin chloride: water is added to the dimethylformamide (DMF) solution in which the reaction occurred, and the product precipitates. Precipitation is useful in purifying many other products: e.g., crude bmim-Cl is taken up in acetonitrile, and dropped into ethyl acetate, where it precipitates.

Biochemistry

Proteins purification and separation can be performed by precipitation in changing the nature of the solvent or the value of its relative permittivity (e.g., by replacing water by ethanol), or by increasing the ionic strength of the solution. As proteins have complex tertiary and quaternary structures due to their specific folding and various weak intermolecular interactions (e.g., hydrogen bridges), these superstructures can be modified and proteins denaturated and precipitated. Another important application of an antisolvent is in ethanol precipitation of DNA.

Metallurgy and alloys

In solid phases, precipitation occurs if the concentration of one solid is above the solubility limit in the host solid, due to e.g. rapid quenching or ion implantation, and the temperature is high enough that diffusion can lead to segregation into precipitates. Precipitation in solids is routinely used to synthesize nanoclusters.

In metallurgy, precipitation from a solid solution is also a way to strengthen alloys.

Precipitation of ceramic phases in metallic alloys such as zirconium hydrides in zircaloy cladding of nuclear fuel pins can also render metallic alloys brittle and lead to their mechanical failure. Correctly mastering the precise temperature and pressure conditions when cooling down spent nuclear fuels is therefore essential to avoid damaging their cladding and to preserve the integrity of the spent fuel elements on the long term in dry storage casks and in geological disposal conditions.

Industrial processes

Hydroxide precipitation is probably the most widely used industrial precipitation process in which metal hydroxides are formed by adding calcium hydroxide (slaked lime) or sodium hydroxide (caustic soda) as precipitant.

History

Powders derived from different precipitation processes have also historically been known as 'flowers'.

Additional Information

Precipitation is the formation of a solid in a solution during a chemical reaction. When the reaction occurs, the solid formed is called the precipitate, and the liquid remaining above the solid is called the supernate.

Uses of precipitation reactions

Precipitation reactions can be used for making pigments, removing salts from water in water treatment, and for qualitative chemical analysis.

This effect is useful in many industrial and scientific applications whereby a chemical reaction may produce a solid that can be collected from the solution by various methods (e.g. filtration, decanting, centrifugation). Precipitation from a solid solution is also a useful way to strengthen alloys; this process is known as solid solution strengthening.

Mechanism

Precipitation can occur when an insoluble substance is formed in the solution due to a chemical reaction or when the solution has been supersaturated by a compound. The formation of a precipitate is a sign of a chemical change. In most situations, the solid forms (“falls”) out of the solute phase, and sink to the bottom of the solution (though it will float if it is less dense than the solvent, or form a suspension).

The solid may reach the bottom of a container by means of settling, sedimentation, or centrifugation.

An important stage of the precipitation process is the onset of nucleation. The creation of a hypothetical solid particle includes the formation of an interface, which requires some energy based on the relative surface energy of the solid and the solution. If this energy is not available, and no suitable nucleation surface is available, supersaturation occurs.

Cation Sensitivity

Precipitate formation is useful in the detection of the type of cation in salt. To do this, an alkali first reacts with the unknown salt to produce a precipitate which is the hydroxide of the unknown salt. To identify the cation, the color of the precipitate and its solubility in excess are noted. Similar processes are often used to separate chemically similar elements, such as the Alkali earth metals.

Digestion

Digestion, or precipitate ageing, happens when a freshly-formed precipitate is left, usually at a higher temperature, in the solution from which it is precipitated. It results in cleaner and bigger particles. The physico-chemical process underlying digestion is called Ostwald ripening.

Coprecipitation

Coprecipitation is the carrying down by a precipitate of substances normally soluble under the conditions employed. It is an important issue in chemical analysis, where it is often undesirable, but in some cases it can be exploited. In gravimetric analysis, it is a problem because undesired impurities often coprecipitate with the analyte, resulting in excess mass. On the other hand, in the analysis of trace elements, as is often the case in radiochemistry, coprecipitation is often the only way of separating an element.

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